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1096 Commits
release-1.
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v2.0.1
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|
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32
.dev_scripts/diff_images.py
Normal file
@@ -0,0 +1,32 @@
|
||||
import argparse
|
||||
|
||||
import numpy as np
|
||||
from PIL import Image
|
||||
|
||||
|
||||
def read_image_int16(image_path):
|
||||
image = Image.open(image_path)
|
||||
return np.array(image).astype(np.int16)
|
||||
|
||||
|
||||
def calc_images_mean_L1(image1_path, image2_path):
|
||||
image1 = read_image_int16(image1_path)
|
||||
image2 = read_image_int16(image2_path)
|
||||
assert image1.shape == image2.shape
|
||||
|
||||
mean_L1 = np.abs(image1 - image2).mean()
|
||||
return mean_L1
|
||||
|
||||
|
||||
def parse_args():
|
||||
parser = argparse.ArgumentParser()
|
||||
parser.add_argument('image1_path')
|
||||
parser.add_argument('image2_path')
|
||||
args = parser.parse_args()
|
||||
return args
|
||||
|
||||
|
||||
if __name__ == '__main__':
|
||||
args = parse_args()
|
||||
mean_L1 = calc_images_mean_L1(args.image1_path, args.image2_path)
|
||||
print(mean_L1)
|
||||
|
After Width: | Height: | Size: 416 KiB |
1
.dev_scripts/sample_command.txt
Normal file
@@ -0,0 +1 @@
|
||||
"a photograph of an astronaut riding a horse" -s50 -S42
|
||||
19
.dev_scripts/test_regression_txt2img_dream_v1_4.sh
Normal file
@@ -0,0 +1,19 @@
|
||||
# generate an image
|
||||
PROMPT_FILE=".dev_scripts/sample_command.txt"
|
||||
OUT_DIR="outputs/img-samples/test_regression_txt2img_v1_4"
|
||||
SAMPLES_DIR=${OUT_DIR}
|
||||
python scripts/dream.py \
|
||||
--from_file ${PROMPT_FILE} \
|
||||
--outdir ${OUT_DIR} \
|
||||
--sampler plms
|
||||
|
||||
# original output by CompVis/stable-diffusion
|
||||
IMAGE1=".dev_scripts/images/v1_4_astronaut_rides_horse_plms_step50_seed42.png"
|
||||
# new output
|
||||
IMAGE2=`ls -A ${SAMPLES_DIR}/*.png | sort | tail -n 1`
|
||||
|
||||
echo ""
|
||||
echo "comparing the following two images"
|
||||
echo "IMAGE1: ${IMAGE1}"
|
||||
echo "IMAGE2: ${IMAGE2}"
|
||||
python .dev_scripts/diff_images.py ${IMAGE1} ${IMAGE2}
|
||||
23
.dev_scripts/test_regression_txt2img_v1_4.sh
Normal file
@@ -0,0 +1,23 @@
|
||||
# generate an image
|
||||
PROMPT="a photograph of an astronaut riding a horse"
|
||||
OUT_DIR="outputs/txt2img-samples/test_regression_txt2img_v1_4"
|
||||
SAMPLES_DIR="outputs/txt2img-samples/test_regression_txt2img_v1_4/samples"
|
||||
python scripts/orig_scripts/txt2img.py \
|
||||
--prompt "${PROMPT}" \
|
||||
--outdir ${OUT_DIR} \
|
||||
--plms \
|
||||
--ddim_steps 50 \
|
||||
--n_samples 1 \
|
||||
--n_iter 1 \
|
||||
--seed 42
|
||||
|
||||
# original output by CompVis/stable-diffusion
|
||||
IMAGE1=".dev_scripts/images/v1_4_astronaut_rides_horse_plms_step50_seed42.png"
|
||||
# new output
|
||||
IMAGE2=`ls -A ${SAMPLES_DIR}/*.png | sort | tail -n 1`
|
||||
|
||||
echo ""
|
||||
echo "comparing the following two images"
|
||||
echo "IMAGE1: ${IMAGE1}"
|
||||
echo "IMAGE2: ${IMAGE2}"
|
||||
python .dev_scripts/diff_images.py ${IMAGE1} ${IMAGE2}
|
||||
4
.gitattributes
vendored
Normal file
@@ -0,0 +1,4 @@
|
||||
# Auto normalizes line endings on commit so devs don't need to change local settings.
|
||||
# Only affects text files and ignores other file types.
|
||||
# For more info see: https://www.aleksandrhovhannisyan.com/blog/crlf-vs-lf-normalizing-line-endings-in-git/
|
||||
* text=auto
|
||||
4
.github/CODEOWNERS
vendored
Normal file
@@ -0,0 +1,4 @@
|
||||
ldm/invoke/pngwriter.py @CapableWeb
|
||||
ldm/invoke/server_legacy.py @CapableWeb
|
||||
scripts/legacy_api.py @CapableWeb
|
||||
tests/legacy_tests.sh @CapableWeb
|
||||
36
.github/ISSUE_TEMPLATE/bug_report.md
vendored
Normal file
@@ -0,0 +1,36 @@
|
||||
---
|
||||
name: Bug report
|
||||
about: Create a report to help us improve
|
||||
title: ''
|
||||
labels: ''
|
||||
assignees: ''
|
||||
|
||||
---
|
||||
|
||||
**Describe your environment**
|
||||
- GPU: [cuda/amd/mps/cpu]
|
||||
- VRAM: [if known]
|
||||
- CPU arch: [x86/arm]
|
||||
- OS: [Linux/Windows/macOS]
|
||||
- Python: [Anaconda/miniconda/miniforge/pyenv/other (explain)]
|
||||
- Branch: [if `git status` says anything other than "On branch main" paste it here]
|
||||
- Commit: [run `git show` and paste the line that starts with "Merge" here]
|
||||
|
||||
**Describe the bug**
|
||||
A clear and concise description of what the bug is.
|
||||
|
||||
**To Reproduce**
|
||||
Steps to reproduce the behavior:
|
||||
1. Go to '...'
|
||||
2. Click on '....'
|
||||
3. Scroll down to '....'
|
||||
4. See error
|
||||
|
||||
**Expected behavior**
|
||||
A clear and concise description of what you expected to happen.
|
||||
|
||||
**Screenshots**
|
||||
If applicable, add screenshots to help explain your problem.
|
||||
|
||||
**Additional context**
|
||||
Add any other context about the problem here.
|
||||
20
.github/ISSUE_TEMPLATE/feature_request.md
vendored
Normal file
@@ -0,0 +1,20 @@
|
||||
---
|
||||
name: Feature request
|
||||
about: Suggest an idea for this project
|
||||
title: ''
|
||||
labels: ''
|
||||
assignees: ''
|
||||
|
||||
---
|
||||
|
||||
**Is your feature request related to a problem? Please describe.**
|
||||
A clear and concise description of what the problem is. Ex. I'm always frustrated when [...]
|
||||
|
||||
**Describe the solution you'd like**
|
||||
A clear and concise description of what you want to happen.
|
||||
|
||||
**Describe alternatives you've considered**
|
||||
A clear and concise description of any alternative solutions or features you've considered.
|
||||
|
||||
**Additional context**
|
||||
Add any other context or screenshots about the feature request here.
|
||||
70
.github/workflows/create-caches.yml
vendored
Normal file
@@ -0,0 +1,70 @@
|
||||
name: Create Caches
|
||||
on:
|
||||
workflow_dispatch
|
||||
jobs:
|
||||
build:
|
||||
strategy:
|
||||
matrix:
|
||||
os: [ ubuntu-latest, macos-12 ]
|
||||
name: Create Caches on ${{ matrix.os }} conda
|
||||
runs-on: ${{ matrix.os }}
|
||||
steps:
|
||||
- name: Set platform variables
|
||||
id: vars
|
||||
run: |
|
||||
if [ "$RUNNER_OS" = "macOS" ]; then
|
||||
echo "::set-output name=ENV_FILE::environment-mac.yml"
|
||||
echo "::set-output name=PYTHON_BIN::/usr/local/miniconda/envs/ldm/bin/python"
|
||||
elif [ "$RUNNER_OS" = "Linux" ]; then
|
||||
echo "::set-output name=ENV_FILE::environment.yml"
|
||||
echo "::set-output name=PYTHON_BIN::/usr/share/miniconda/envs/ldm/bin/python"
|
||||
fi
|
||||
- name: Checkout sources
|
||||
uses: actions/checkout@v3
|
||||
- name: Use Cached Stable Diffusion v1.4 Model
|
||||
id: cache-sd-v1-4
|
||||
uses: actions/cache@v3
|
||||
env:
|
||||
cache-name: cache-sd-v1-4
|
||||
with:
|
||||
path: models/ldm/stable-diffusion-v1/model.ckpt
|
||||
key: ${{ env.cache-name }}
|
||||
restore-keys: |
|
||||
${{ env.cache-name }}
|
||||
- name: Download Stable Diffusion v1.4 Model
|
||||
if: ${{ steps.cache-sd-v1-4.outputs.cache-hit != 'true' }}
|
||||
run: |
|
||||
if [ ! -e models/ldm/stable-diffusion-v1 ]; then
|
||||
mkdir -p models/ldm/stable-diffusion-v1
|
||||
fi
|
||||
if [ ! -e models/ldm/stable-diffusion-v1/model.ckpt ]; then
|
||||
curl -o models/ldm/stable-diffusion-v1/model.ckpt ${{ secrets.SD_V1_4_URL }}
|
||||
fi
|
||||
- name: Use Cached Dependencies
|
||||
id: cache-conda-env-ldm
|
||||
uses: actions/cache@v3
|
||||
env:
|
||||
cache-name: cache-conda-env-ldm
|
||||
with:
|
||||
path: ~/.conda/envs/ldm
|
||||
key: ${{ env.cache-name }}
|
||||
restore-keys: |
|
||||
${{ env.cache-name }}-${{ runner.os }}-${{ hashFiles(steps.vars.outputs.ENV_FILE) }}
|
||||
- name: Install Dependencies
|
||||
if: ${{ steps.cache-conda-env-ldm.outputs.cache-hit != 'true' }}
|
||||
run: |
|
||||
conda env create -f ${{ steps.vars.outputs.ENV_FILE }}
|
||||
- name: Use Cached Huggingface and Torch models
|
||||
id: cache-huggingface-torch
|
||||
uses: actions/cache@v3
|
||||
env:
|
||||
cache-name: cache-huggingface-torch
|
||||
with:
|
||||
path: ~/.cache
|
||||
key: ${{ env.cache-name }}
|
||||
restore-keys: |
|
||||
${{ env.cache-name }}-${{ hashFiles('scripts/preload_models.py') }}
|
||||
- name: Download Huggingface and Torch models
|
||||
if: ${{ steps.cache-huggingface-torch.outputs.cache-hit != 'true' }}
|
||||
run: |
|
||||
${{ steps.vars.outputs.PYTHON_BIN }} scripts/preload_models.py
|
||||
28
.github/workflows/mkdocs-flow.yml
vendored
Normal file
@@ -0,0 +1,28 @@
|
||||
name: Deploy
|
||||
on:
|
||||
push:
|
||||
branches:
|
||||
- main
|
||||
# pull_request:
|
||||
# branches:
|
||||
# - main
|
||||
jobs:
|
||||
build:
|
||||
name: Deploy docs to GitHub Pages
|
||||
runs-on: ubuntu-latest
|
||||
steps:
|
||||
- name: Checkout
|
||||
uses: actions/checkout@v3
|
||||
with:
|
||||
fetch-depth: 0
|
||||
- name: Build
|
||||
uses: Tiryoh/actions-mkdocs@v0
|
||||
with:
|
||||
mkdocs_version: 'latest' # option
|
||||
requirements: '/requirements-mkdocs.txt' # option
|
||||
configfile: '/mkdocs.yml' # option
|
||||
- name: Deploy
|
||||
uses: peaceiris/actions-gh-pages@v3
|
||||
with:
|
||||
github_token: ${{ secrets.GITHUB_TOKEN }}
|
||||
publish_dir: ./site
|
||||
97
.github/workflows/test-invoke-conda.yml
vendored
Normal file
@@ -0,0 +1,97 @@
|
||||
name: Test Invoke with Conda
|
||||
on:
|
||||
push:
|
||||
branches:
|
||||
- 'main'
|
||||
- 'development'
|
||||
jobs:
|
||||
os_matrix:
|
||||
strategy:
|
||||
matrix:
|
||||
os: [ ubuntu-latest, macos-12 ]
|
||||
name: Test invoke.py on ${{ matrix.os }} with conda
|
||||
runs-on: ${{ matrix.os }}
|
||||
steps:
|
||||
- run: |
|
||||
echo The PR was merged
|
||||
- name: Set platform variables
|
||||
id: vars
|
||||
run: |
|
||||
# Note, can't "activate" via github action; specifying the env's python has the same effect
|
||||
if [ "$RUNNER_OS" = "macOS" ]; then
|
||||
echo "::set-output name=ENV_FILE::environment-mac.yml"
|
||||
echo "::set-output name=PYTHON_BIN::/usr/local/miniconda/envs/ldm/bin/python"
|
||||
elif [ "$RUNNER_OS" = "Linux" ]; then
|
||||
echo "::set-output name=ENV_FILE::environment.yml"
|
||||
echo "::set-output name=PYTHON_BIN::/usr/share/miniconda/envs/ldm/bin/python"
|
||||
fi
|
||||
- name: Checkout sources
|
||||
uses: actions/checkout@v3
|
||||
- name: Use Cached Stable Diffusion v1.4 Model
|
||||
id: cache-sd-v1-4
|
||||
uses: actions/cache@v3
|
||||
env:
|
||||
cache-name: cache-sd-v1-4
|
||||
with:
|
||||
path: models/ldm/stable-diffusion-v1/model.ckpt
|
||||
key: ${{ env.cache-name }}
|
||||
restore-keys: |
|
||||
${{ env.cache-name }}
|
||||
- name: Download Stable Diffusion v1.4 Model
|
||||
if: ${{ steps.cache-sd-v1-4.outputs.cache-hit != 'true' }}
|
||||
run: |
|
||||
if [ ! -e models/ldm/stable-diffusion-v1 ]; then
|
||||
mkdir -p models/ldm/stable-diffusion-v1
|
||||
fi
|
||||
if [ ! -e models/ldm/stable-diffusion-v1/model.ckpt ]; then
|
||||
curl -o models/ldm/stable-diffusion-v1/model.ckpt ${{ secrets.SD_V1_4_URL }}
|
||||
fi
|
||||
- name: Use Cached Dependencies
|
||||
id: cache-conda-env-ldm
|
||||
uses: actions/cache@v3
|
||||
env:
|
||||
cache-name: cache-conda-env-ldm
|
||||
with:
|
||||
path: ~/.conda/envs/ldm
|
||||
key: ${{ env.cache-name }}
|
||||
restore-keys: |
|
||||
${{ env.cache-name }}-${{ runner.os }}-${{ hashFiles(steps.vars.outputs.ENV_FILE) }}
|
||||
- name: Install Dependencies
|
||||
if: ${{ steps.cache-conda-env-ldm.outputs.cache-hit != 'true' }}
|
||||
run: |
|
||||
conda env create -f ${{ steps.vars.outputs.ENV_FILE }}
|
||||
- name: Use Cached Huggingface and Torch models
|
||||
id: cache-hugginface-torch
|
||||
uses: actions/cache@v3
|
||||
env:
|
||||
cache-name: cache-hugginface-torch
|
||||
with:
|
||||
path: ~/.cache
|
||||
key: ${{ env.cache-name }}
|
||||
restore-keys: |
|
||||
${{ env.cache-name }}-${{ hashFiles('scripts/preload_models.py') }}
|
||||
- name: Download Huggingface and Torch models
|
||||
if: ${{ steps.cache-hugginface-torch.outputs.cache-hit != 'true' }}
|
||||
run: |
|
||||
${{ steps.vars.outputs.PYTHON_BIN }} scripts/preload_models.py
|
||||
# - name: Run tmate
|
||||
# uses: mxschmitt/action-tmate@v3
|
||||
# timeout-minutes: 30
|
||||
- name: Run the tests
|
||||
run: |
|
||||
# Note, can't "activate" via github action; specifying the env's python has the same effect
|
||||
if [ $(uname) = "Darwin" ]; then
|
||||
export PYTORCH_ENABLE_MPS_FALLBACK=1
|
||||
fi
|
||||
# Utterly hacky, but I don't know how else to do this
|
||||
if [[ ${{ github.ref }} == 'refs/heads/master' ]]; then
|
||||
time ${{ steps.vars.outputs.PYTHON_BIN }} scripts/invoke.py --from_file tests/preflight_prompts.txt
|
||||
elif [[ ${{ github.ref }} == 'refs/heads/development' ]]; then
|
||||
time ${{ steps.vars.outputs.PYTHON_BIN }} scripts/invoke.py --from_file tests/dev_prompts.txt
|
||||
fi
|
||||
mkdir -p outputs/img-samples
|
||||
- name: Archive results
|
||||
uses: actions/upload-artifact@v3
|
||||
with:
|
||||
name: results
|
||||
path: outputs/img-samples
|
||||
28
.gitignore
vendored
@@ -1,6 +1,13 @@
|
||||
# ignore default image save location and model symbolic link
|
||||
outputs/
|
||||
models/ldm/stable-diffusion-v1/model.ckpt
|
||||
ldm/dream/restoration/codeformer/weights
|
||||
|
||||
# ignore the Anaconda/Miniconda installer used while building Docker image
|
||||
anaconda.sh
|
||||
|
||||
# ignore a directory which serves as a place for initial images
|
||||
inputs/
|
||||
|
||||
# Byte-compiled / optimized / DLL files
|
||||
__pycache__/
|
||||
@@ -170,6 +177,25 @@ cython_debug/
|
||||
#.idea/
|
||||
|
||||
src
|
||||
logs/
|
||||
**/__pycache__/
|
||||
outputs
|
||||
|
||||
# Logs and associated folders
|
||||
# created from generated embeddings.
|
||||
logs
|
||||
testtube
|
||||
checkpoints
|
||||
# If it's a Mac
|
||||
.DS_Store
|
||||
|
||||
# Let the frontend manage its own gitignore
|
||||
!frontend/*
|
||||
|
||||
# Scratch folder
|
||||
.scratch/
|
||||
.vscode/
|
||||
gfpgan/
|
||||
models/ldm/stable-diffusion-v1/model.sha256
|
||||
|
||||
# GFPGAN model files
|
||||
gfpgan/
|
||||
|
||||
13
.prettierrc.yaml
Normal file
@@ -0,0 +1,13 @@
|
||||
endOfLine: lf
|
||||
tabWidth: 2
|
||||
useTabs: false
|
||||
singleQuote: true
|
||||
quoteProps: as-needed
|
||||
embeddedLanguageFormatting: auto
|
||||
overrides:
|
||||
- files: '*.md'
|
||||
options:
|
||||
proseWrap: always
|
||||
printWidth: 80
|
||||
parser: markdown
|
||||
cursorOffset: -1
|
||||
@@ -1,210 +0,0 @@
|
||||
# Original README from CompViz/stable-diffusion
|
||||
*Stable Diffusion was made possible thanks to a collaboration with [Stability AI](https://stability.ai/) and [Runway](https://runwayml.com/) and builds upon our previous work:*
|
||||
|
||||
[**High-Resolution Image Synthesis with Latent Diffusion Models**](https://ommer-lab.com/research/latent-diffusion-models/)<br/>
|
||||
[Robin Rombach](https://github.com/rromb)\*,
|
||||
[Andreas Blattmann](https://github.com/ablattmann)\*,
|
||||
[Dominik Lorenz](https://github.com/qp-qp)\,
|
||||
[Patrick Esser](https://github.com/pesser),
|
||||
[Björn Ommer](https://hci.iwr.uni-heidelberg.de/Staff/bommer)<br/>
|
||||
|
||||
**CVPR '22 Oral**
|
||||
|
||||
which is available on [GitHub](https://github.com/CompVis/latent-diffusion). PDF at [arXiv](https://arxiv.org/abs/2112.10752). Please also visit our [Project page](https://ommer-lab.com/research/latent-diffusion-models/).
|
||||
|
||||

|
||||
[Stable Diffusion](#stable-diffusion-v1) is a latent text-to-image diffusion
|
||||
model.
|
||||
Thanks to a generous compute donation from [Stability AI](https://stability.ai/) and support from [LAION](https://laion.ai/), we were able to train a Latent Diffusion Model on 512x512 images from a subset of the [LAION-5B](https://laion.ai/blog/laion-5b/) database.
|
||||
Similar to Google's [Imagen](https://arxiv.org/abs/2205.11487),
|
||||
this model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts.
|
||||
With its 860M UNet and 123M text encoder, the model is relatively lightweight and runs on a GPU with at least 10GB VRAM.
|
||||
See [this section](#stable-diffusion-v1) below and the [model card](https://huggingface.co/CompVis/stable-diffusion).
|
||||
|
||||
|
||||
## Requirements
|
||||
|
||||
A suitable [conda](https://conda.io/) environment named `ldm` can be created
|
||||
and activated with:
|
||||
|
||||
```
|
||||
conda env create -f environment.yaml
|
||||
conda activate ldm
|
||||
```
|
||||
|
||||
You can also update an existing [latent diffusion](https://github.com/CompVis/latent-diffusion) environment by running
|
||||
|
||||
```
|
||||
conda install pytorch torchvision -c pytorch
|
||||
pip install transformers==4.19.2
|
||||
pip install -e .
|
||||
```
|
||||
|
||||
## Stable Diffusion v1
|
||||
|
||||
Stable Diffusion v1 refers to a specific configuration of the model
|
||||
architecture that uses a downsampling-factor 8 autoencoder with an 860M UNet
|
||||
and CLIP ViT-L/14 text encoder for the diffusion model. The model was pretrained on 256x256 images and
|
||||
then finetuned on 512x512 images.
|
||||
|
||||
*Note: Stable Diffusion v1 is a general text-to-image diffusion model and therefore mirrors biases and (mis-)conceptions that are present
|
||||
in its training data.
|
||||
Details on the training procedure and data, as well as the intended use of the model can be found in the corresponding [model card](https://huggingface.co/CompVis/stable-diffusion).
|
||||
Research into the safe deployment of general text-to-image models is an ongoing effort. To prevent misuse and harm, we currently provide access to the checkpoints only for [academic research purposes upon request](https://stability.ai/academia-access-form).
|
||||
**This is an experiment in safe and community-driven publication of a capable and general text-to-image model. We are working on a public release with a more permissive license that also incorporates ethical considerations.***
|
||||
|
||||
[Request access to Stable Diffusion v1 checkpoints for academic research](https://stability.ai/academia-access-form)
|
||||
|
||||
### Weights
|
||||
|
||||
We currently provide three checkpoints, `sd-v1-1.ckpt`, `sd-v1-2.ckpt` and `sd-v1-3.ckpt`,
|
||||
which were trained as follows,
|
||||
|
||||
- `sd-v1-1.ckpt`: 237k steps at resolution `256x256` on [laion2B-en](https://huggingface.co/datasets/laion/laion2B-en).
|
||||
194k steps at resolution `512x512` on [laion-high-resolution](https://huggingface.co/datasets/laion/laion-high-resolution) (170M examples from LAION-5B with resolution `>= 1024x1024`).
|
||||
- `sd-v1-2.ckpt`: Resumed from `sd-v1-1.ckpt`.
|
||||
515k steps at resolution `512x512` on "laion-improved-aesthetics" (a subset of laion2B-en,
|
||||
filtered to images with an original size `>= 512x512`, estimated aesthetics score `> 5.0`, and an estimated watermark probability `< 0.5`. The watermark estimate is from the LAION-5B metadata, the aesthetics score is estimated using an [improved aesthetics estimator](https://github.com/christophschuhmann/improved-aesthetic-predictor)).
|
||||
- `sd-v1-3.ckpt`: Resumed from `sd-v1-2.ckpt`. 195k steps at resolution `512x512` on "laion-improved-aesthetics" and 10\% dropping of the text-conditioning to improve [classifier-free guidance sampling](https://arxiv.org/abs/2207.12598).
|
||||
|
||||
Evaluations with different classifier-free guidance scales (1.5, 2.0, 3.0, 4.0,
|
||||
5.0, 6.0, 7.0, 8.0) and 50 PLMS sampling
|
||||
steps show the relative improvements of the checkpoints:
|
||||

|
||||
|
||||
|
||||
|
||||
### Text-to-Image with Stable Diffusion
|
||||

|
||||

|
||||
|
||||
Stable Diffusion is a latent diffusion model conditioned on the (non-pooled) text embeddings of a CLIP ViT-L/14 text encoder.
|
||||
|
||||
|
||||
#### Sampling Script
|
||||
|
||||
After [obtaining the weights](#weights), link them
|
||||
```
|
||||
mkdir -p models/ldm/stable-diffusion-v1/
|
||||
ln -s <path/to/model.ckpt> models/ldm/stable-diffusion-v1/model.ckpt
|
||||
```
|
||||
and sample with
|
||||
```
|
||||
python scripts/txt2img.py --prompt "a photograph of an astronaut riding a horse" --plms
|
||||
```
|
||||
|
||||
By default, this uses a guidance scale of `--scale 7.5`, [Katherine Crowson's implementation](https://github.com/CompVis/latent-diffusion/pull/51) of the [PLMS](https://arxiv.org/abs/2202.09778) sampler,
|
||||
and renders images of size 512x512 (which it was trained on) in 50 steps. All supported arguments are listed below (type `python scripts/txt2img.py --help`).
|
||||
|
||||
```commandline
|
||||
usage: txt2img.py [-h] [--prompt [PROMPT]] [--outdir [OUTDIR]] [--skip_grid] [--skip_save] [--ddim_steps DDIM_STEPS] [--plms] [--laion400m] [--fixed_code] [--ddim_eta DDIM_ETA] [--n_iter N_ITER] [--H H] [--W W] [--C C] [--f F] [--n_samples N_SAMPLES] [--n_rows N_ROWS]
|
||||
[--scale SCALE] [--from-file FROM_FILE] [--config CONFIG] [--ckpt CKPT] [--seed SEED] [--precision {full,autocast}]
|
||||
|
||||
optional arguments:
|
||||
-h, --help show this help message and exit
|
||||
--prompt [PROMPT] the prompt to render
|
||||
--outdir [OUTDIR] dir to write results to
|
||||
--skip_grid do not save a grid, only individual samples. Helpful when evaluating lots of samples
|
||||
--skip_save do not save individual samples. For speed measurements.
|
||||
--ddim_steps DDIM_STEPS
|
||||
number of ddim sampling steps
|
||||
--plms use plms sampling
|
||||
--laion400m uses the LAION400M model
|
||||
--fixed_code if enabled, uses the same starting code across samples
|
||||
--ddim_eta DDIM_ETA ddim eta (eta=0.0 corresponds to deterministic sampling
|
||||
--n_iter N_ITER sample this often
|
||||
--H H image height, in pixel space
|
||||
--W W image width, in pixel space
|
||||
--C C latent channels
|
||||
--f F downsampling factor
|
||||
--n_samples N_SAMPLES
|
||||
how many samples to produce for each given prompt. A.k.a. batch size
|
||||
(note that the seeds for each image in the batch will be unavailable)
|
||||
--n_rows N_ROWS rows in the grid (default: n_samples)
|
||||
--scale SCALE unconditional guidance scale: eps = eps(x, empty) + scale * (eps(x, cond) - eps(x, empty))
|
||||
--from-file FROM_FILE
|
||||
if specified, load prompts from this file
|
||||
--config CONFIG path to config which constructs model
|
||||
--ckpt CKPT path to checkpoint of model
|
||||
--seed SEED the seed (for reproducible sampling)
|
||||
--precision {full,autocast}
|
||||
evaluate at this precision
|
||||
|
||||
```
|
||||
Note: The inference config for all v1 versions is designed to be used with EMA-only checkpoints.
|
||||
For this reason `use_ema=False` is set in the configuration, otherwise the code will try to switch from
|
||||
non-EMA to EMA weights. If you want to examine the effect of EMA vs no EMA, we provide "full" checkpoints
|
||||
which contain both types of weights. For these, `use_ema=False` will load and use the non-EMA weights.
|
||||
|
||||
|
||||
#### Diffusers Integration
|
||||
|
||||
Another way to download and sample Stable Diffusion is by using the [diffusers library](https://github.com/huggingface/diffusers/tree/main#new--stable-diffusion-is-now-fully-compatible-with-diffusers)
|
||||
```py
|
||||
# make sure you're logged in with `huggingface-cli login`
|
||||
from torch import autocast
|
||||
from diffusers import StableDiffusionPipeline, LMSDiscreteScheduler
|
||||
|
||||
pipe = StableDiffusionPipeline.from_pretrained(
|
||||
"CompVis/stable-diffusion-v1-3-diffusers",
|
||||
use_auth_token=True
|
||||
)
|
||||
|
||||
prompt = "a photo of an astronaut riding a horse on mars"
|
||||
with autocast("cuda"):
|
||||
image = pipe(prompt)["sample"][0]
|
||||
|
||||
image.save("astronaut_rides_horse.png")
|
||||
```
|
||||
|
||||
|
||||
|
||||
### Image Modification with Stable Diffusion
|
||||
|
||||
By using a diffusion-denoising mechanism as first proposed by [SDEdit](https://arxiv.org/abs/2108.01073), the model can be used for different
|
||||
tasks such as text-guided image-to-image translation and upscaling. Similar to the txt2img sampling script,
|
||||
we provide a script to perform image modification with Stable Diffusion.
|
||||
|
||||
The following describes an example where a rough sketch made in [Pinta](https://www.pinta-project.com/) is converted into a detailed artwork.
|
||||
```
|
||||
python scripts/img2img.py --prompt "A fantasy landscape, trending on artstation" --init-img <path-to-img.jpg> --strength 0.8
|
||||
```
|
||||
Here, strength is a value between 0.0 and 1.0, that controls the amount of noise that is added to the input image.
|
||||
Values that approach 1.0 allow for lots of variations but will also produce images that are not semantically consistent with the input. See the following example.
|
||||
|
||||
**Input**
|
||||
|
||||

|
||||
|
||||
**Outputs**
|
||||
|
||||

|
||||

|
||||
|
||||
This procedure can, for example, also be used to upscale samples from the base model.
|
||||
|
||||
|
||||
## Comments
|
||||
|
||||
- Our codebase for the diffusion models builds heavily on [OpenAI's ADM codebase](https://github.com/openai/guided-diffusion)
|
||||
and [https://github.com/lucidrains/denoising-diffusion-pytorch](https://github.com/lucidrains/denoising-diffusion-pytorch).
|
||||
Thanks for open-sourcing!
|
||||
|
||||
- The implementation of the transformer encoder is from [x-transformers](https://github.com/lucidrains/x-transformers) by [lucidrains](https://github.com/lucidrains?tab=repositories).
|
||||
|
||||
|
||||
## BibTeX
|
||||
|
||||
```
|
||||
@misc{rombach2021highresolution,
|
||||
title={High-Resolution Image Synthesis with Latent Diffusion Models},
|
||||
author={Robin Rombach and Andreas Blattmann and Dominik Lorenz and Patrick Esser and Björn Ommer},
|
||||
year={2021},
|
||||
eprint={2112.10752},
|
||||
archivePrefix={arXiv},
|
||||
primaryClass={cs.CV}
|
||||
}
|
||||
|
||||
```
|
||||
|
||||
|
||||
602
README.md
@@ -1,483 +1,205 @@
|
||||
# Stable Diffusion Dream Script
|
||||
<div align="center">
|
||||
|
||||
This is a fork of CompVis/stable-diffusion, the wonderful open source
|
||||
text-to-image generator.
|
||||
# InvokeAI: A Stable Diffusion Toolkit
|
||||
|
||||
The original has been modified in several ways:
|
||||
_Formally known as lstein/stable-diffusion_
|
||||
|
||||
## Interactive command-line interface similar to the Discord bot
|
||||

|
||||
|
||||
The *dream.py* script, located in scripts/dream.py,
|
||||
provides an interactive interface to image generation similar to
|
||||
the "dream mothership" bot that Stable AI provided on its Discord
|
||||
server. Unlike the txt2img.py and img2img.py scripts provided in the
|
||||
original CompViz/stable-diffusion source code repository, the
|
||||
time-consuming initialization of the AI model
|
||||
initialization only happens once. After that image generation
|
||||
from the command-line interface is very fast.
|
||||
[![discord badge]][discord link]
|
||||
|
||||
The script uses the readline library to allow for in-line editing,
|
||||
command history (up and down arrows), autocompletion, and more. To help
|
||||
keep track of which prompts generated which images, the script writes a
|
||||
log file of image names and prompts to the selected output directory.
|
||||
In addition, as of version 1.02, it also writes the prompt into the PNG
|
||||
file's metadata where it can be retrieved using scripts/images2prompt.py
|
||||
[![latest release badge]][latest release link] [![github stars badge]][github stars link] [![github forks badge]][github forks link]
|
||||
|
||||
The script is confirmed to work on Linux and Windows systems. It should
|
||||
work on MacOSX as well, but this is not confirmed. Note that this script
|
||||
runs from the command-line (CMD or Terminal window), and does not have a GUI.
|
||||
[![CI checks on main badge]][CI checks on main link] [![CI checks on dev badge]][CI checks on dev link] [![latest commit to dev badge]][latest commit to dev link]
|
||||
|
||||
~~~~
|
||||
(ldm) ~/stable-diffusion$ python3 ./scripts/dream.py
|
||||
* Initializing, be patient...
|
||||
Loading model from models/ldm/text2img-large/model.ckpt
|
||||
(...more initialization messages...)
|
||||
[![github open issues badge]][github open issues link] [![github open prs badge]][github open prs link]
|
||||
|
||||
* Initialization done! Awaiting your command...
|
||||
dream> ashley judd riding a camel -n2 -s150
|
||||
Outputs:
|
||||
outputs/img-samples/00009.png: "ashley judd riding a camel" -n2 -s150 -S 416354203
|
||||
outputs/img-samples/00010.png: "ashley judd riding a camel" -n2 -s150 -S 1362479620
|
||||
[CI checks on dev badge]: https://flat.badgen.net/github/checks/invoke-ai/InvokeAI/development?label=CI%20status%20on%20dev&cache=900&icon=github
|
||||
[CI checks on dev link]: https://github.com/invoke-ai/InvokeAI/actions?query=branch%3Adevelopment
|
||||
[CI checks on main badge]: https://flat.badgen.net/github/checks/invoke-ai/InvokeAI/main?label=CI%20status%20on%20main&cache=900&icon=github
|
||||
[CI checks on main link]: https://github.com/invoke-ai/InvokeAI/actions/workflows/test-invoke-conda.yml
|
||||
[discord badge]: https://flat.badgen.net/discord/members/ZmtBAhwWhy?icon=discord
|
||||
[discord link]: https://discord.gg/ZmtBAhwWhy
|
||||
[github forks badge]: https://flat.badgen.net/github/forks/invoke-ai/InvokeAI?icon=github
|
||||
[github forks link]: https://useful-forks.github.io/?repo=invoke-ai%2FInvokeAI
|
||||
[github open issues badge]: https://flat.badgen.net/github/open-issues/invoke-ai/InvokeAI?icon=github
|
||||
[github open issues link]: https://github.com/invoke-ai/InvokeAI/issues?q=is%3Aissue+is%3Aopen
|
||||
[github open prs badge]: https://flat.badgen.net/github/open-prs/invoke-ai/InvokeAI?icon=github
|
||||
[github open prs link]: https://github.com/invoke-ai/InvokeAI/pulls?q=is%3Apr+is%3Aopen
|
||||
[github stars badge]: https://flat.badgen.net/github/stars/invoke-ai/InvokeAI?icon=github
|
||||
[github stars link]: https://github.com/invoke-ai/InvokeAI/stargazers
|
||||
[latest commit to dev badge]: https://flat.badgen.net/github/last-commit/invoke-ai/InvokeAI/development?icon=github&color=yellow&label=last%20dev%20commit&cache=900
|
||||
[latest commit to dev link]: https://github.com/invoke-ai/InvokeAI/commits/development
|
||||
[latest release badge]: https://flat.badgen.net/github/release/invoke-ai/InvokeAI/development?icon=github
|
||||
[latest release link]: https://github.com/invoke-ai/InvokeAI/releases
|
||||
</div>
|
||||
|
||||
dream> "there's a fly in my soup" -n6 -g
|
||||
outputs/img-samples/00011.png: "there's a fly in my soup" -n6 -g -S 2685670268
|
||||
seeds for individual rows: [2685670268, 1216708065, 2335773498, 822223658, 714542046, 3395302430]
|
||||
dream> q
|
||||
This is a fork of
|
||||
[CompVis/stable-diffusion](https://github.com/CompVis/stable-diffusion),
|
||||
the open source text-to-image generator. It provides a streamlined
|
||||
process with various new features and options to aid the image
|
||||
generation process. It runs on Windows, Mac and Linux machines, with
|
||||
GPU cards with as little as 4 GB of RAM. It provides both a polished
|
||||
Web interface (see below), and an easy-to-use command-line interface.
|
||||
|
||||
# this shows how to retrieve the prompt stored in the saved image's metadata
|
||||
(ldm) ~/stable-diffusion$ python3 ./scripts/images2prompt.py outputs/img_samples/*.png
|
||||
00009.png: "ashley judd riding a camel" -s150 -S 416354203
|
||||
00010.png: "ashley judd riding a camel" -s150 -S 1362479620
|
||||
00011.png: "there's a fly in my soup" -n6 -g -S 2685670268
|
||||
~~~~
|
||||
**Quick links**: [<a href="https://discord.gg/NwVCmKwY">Discord Server</a>] [<a href="https://invoke-ai.github.io/InvokeAI/">Documentation and Tutorials</a>] [<a href="https://github.com/invoke-ai/InvokeAI/">Code and Downloads</a>] [<a href="https://github.com/invoke-ai/InvokeAI/issues">Bug Reports</a>] [<a href="https://github.com/invoke-ai/InvokeAI/discussions">Discussion, Ideas & Q&A</a>]
|
||||
|
||||
The dream> prompt's arguments are pretty much identical to those used
|
||||
in the Discord bot, except you don't need to type "!dream" (it doesn't
|
||||
hurt if you do). A significant change is that creation of individual
|
||||
images is now the default unless --grid (-g) is given. For backward
|
||||
compatibility, the -i switch is recognized. For command-line help
|
||||
type -h (or --help) at the dream> prompt.
|
||||
|
||||
The script itself also recognizes a series of command-line switches
|
||||
that will change important global defaults, such as the directory for
|
||||
image outputs and the location of the model weight files.
|
||||
|
||||
## Image-to-Image
|
||||
|
||||
This script also provides an img2img feature that lets you seed your
|
||||
creations with a drawing or photo. This is a really cool feature that tells
|
||||
stable diffusion to build the prompt on top of the image you provide, preserving
|
||||
the original's basic shape and layout. To use it, provide the --init_img
|
||||
option as shown here:
|
||||
|
||||
~~~~
|
||||
dream> "waterfall and rainbow" --init_img=./init-images/crude_drawing.png --strength=0.5 -s100 -n4
|
||||
~~~~
|
||||
|
||||
The --init_img (-I) option gives the path to the seed picture. --strength (-f) controls how much
|
||||
the original will be modified, ranging from 0.0 (keep the original intact), to 1.0 (ignore the original
|
||||
completely). The default is 0.75, and ranges from 0.25-0.75 give interesting results.
|
||||
|
||||
You may also pass a -v<count> option to generate count variants on the original image. This is done by
|
||||
passing the first generated image back into img2img the requested number of times. It generates interesting
|
||||
variants.
|
||||
|
||||
## Barebones Web Server
|
||||
|
||||
As of version 1.10, this distribution comes with a bare bones web server (see screenshot). To use it,
|
||||
run the command:
|
||||
|
||||
~~~~
|
||||
(ldm) ~/stable-diffusion$ python3 scripts/dream_web.py
|
||||
~~~~
|
||||
|
||||
You can then connect to the server by pointing your web browser at
|
||||
http://localhost:9090, or to the network name or IP address of the server.
|
||||
|
||||
Kudos to [Tesseract Cat](https://github.com/TesseractCat) for
|
||||
contributing this code.
|
||||
|
||||

|
||||
|
||||
## Weighted Prompts
|
||||
|
||||
You may weight different sections of the prompt to tell the sampler to attach different levels of
|
||||
priority to them, by adding :(number) to the end of the section you wish to up- or downweight.
|
||||
For example consider this prompt:
|
||||
|
||||
~~~~
|
||||
tabby cat:0.25 white duck:0.75 hybrid
|
||||
~~~~
|
||||
|
||||
This will tell the sampler to invest 25% of its effort on the tabby
|
||||
cat aspect of the image and 75% on the white duck aspect
|
||||
(surprisingly, this example actually works). The prompt weights can
|
||||
use any combination of integers and floating point numbers, and they
|
||||
do not need to add up to 1.
|
||||
|
||||
## Personalizing Text-to-Image Generation
|
||||
|
||||
You may personalize the generated images to provide your own styles or objects by training a new LDM checkpoint
|
||||
and introducing a new vocabulary to the fixed model.
|
||||
|
||||
To train, prepare a folder that contains images sized at 512x512 and execute the following:
|
||||
|
||||
~~~~
|
||||
# As the default backend is not available on Windows, if you're using that platform, execute SET PL_TORCH_DISTRIBUTED_BACKEND=gloo
|
||||
(ldm) ~/stable-diffusion$ python3 ./main.py --base ./configs/stable-diffusion/v1-finetune.yaml \
|
||||
-t \
|
||||
--actual_resume ./models/ldm/stable-diffusion-v1/model.ckpt \
|
||||
-n my_cat \
|
||||
--gpus 0, \
|
||||
--data_root D:/textual-inversion/my_cat \
|
||||
--init_word 'cat'
|
||||
~~~~
|
||||
|
||||
During the training process, files will be created in /logs/[project][time][project]/
|
||||
where you can see the process.
|
||||
|
||||
conditioning* contains the training prompts
|
||||
inputs, reconstruction the input images for the training epoch
|
||||
samples, samples scaled for a sample of the prompt and one with the init word provided
|
||||
|
||||
On a RTX3090, the process for SD will take ~1h @1.6 iterations/sec.
|
||||
|
||||
Note: According to the associated paper, the optimal number of images
|
||||
is 3-5. Your model may not converge if you use more images than that.
|
||||
|
||||
Training will run indefinately, but you may wish to stop it before the
|
||||
heat death of the universe, when you find a low loss epoch or around
|
||||
~5000 iterations.
|
||||
|
||||
Once the model is trained, specify the trained .pt file when starting
|
||||
dream using
|
||||
|
||||
~~~~
|
||||
(ldm) ~/stable-diffusion$ python3 ./scripts/dream.py --embedding_path /path/to/embedding.pt --full_precision
|
||||
~~~~
|
||||
|
||||
Then, to utilize your subject at the dream prompt
|
||||
|
||||
~~~
|
||||
dream> "a photo of *"
|
||||
~~~
|
||||
|
||||
this also works with image2image
|
||||
~~~~
|
||||
dream> "waterfall and rainbow in the style of *" --init_img=./init-images/crude_drawing.png --strength=0.5 -s100 -n4
|
||||
~~~~
|
||||
|
||||
It's also possible to train multiple tokens (modify the placeholder string in configs/stable-diffusion/v1-finetune.yaml) and combine LDM checkpoints using:
|
||||
|
||||
~~~~
|
||||
(ldm) ~/stable-diffusion$ python3 ./scripts/merge_embeddings.py \
|
||||
--manager_ckpts /path/to/first/embedding.pt /path/to/second/embedding.pt [...] \
|
||||
--output_path /path/to/output/embedding.pt
|
||||
~~~~
|
||||
|
||||
Credit goes to @rinongal and the repository located at
|
||||
https://github.com/rinongal/textual_inversion Please see the
|
||||
repository and associated paper for details and limitations.
|
||||
|
||||
## Changes
|
||||
|
||||
* v1.09 (24 August 2022)
|
||||
* A barebone web server for interactive online generation of txt2img and img2img.
|
||||
* A new -v option allows you to generate multiple variants of an initial image
|
||||
in img2img mode. (kudos to [Oceanswave](https://github.com/Oceanswave). [See this discussion in the PR for examples and details on use](https://github.com/lstein/stable-diffusion/pull/71#issuecomment-1226700810))
|
||||
* Added ability to personalize text to image generation (kudos to [Oceanswave](https://github.com/Oceanswave) and [nicolai256](https://github.com/nicolai256))
|
||||
* Enabled all of the samplers from k_diffusion
|
||||
|
||||
* v1.08 (24 August 2022)
|
||||
* Escape single quotes on the dream> command before trying to parse. This avoids
|
||||
parse errors.
|
||||
* Removed instruction to get Python3.8 as first step in Windows install.
|
||||
Anaconda3 does it for you.
|
||||
* Added bounds checks for numeric arguments that could cause crashes.
|
||||
* Cleaned up the copyright and license agreement files.
|
||||
|
||||
* v1.07 (23 August 2022)
|
||||
* Image filenames will now never fill gaps in the sequence, but will be assigned the
|
||||
next higher name in the chosen directory. This ensures that the alphabetic and chronological
|
||||
sort orders are the same.
|
||||
|
||||
* v1.06 (23 August 2022)
|
||||
* Added weighted prompt support contributed by [xraxra](https://github.com/xraxra)
|
||||
* Example of using weighted prompts to tweak a demonic figure contributed by [bmaltais](https://github.com/bmaltais)
|
||||
|
||||
* v1.05 (22 August 2022 - after the drop)
|
||||
* Filenames now use the following formats:
|
||||
000010.95183149.png -- Two files produced by the same command (e.g. -n2),
|
||||
000010.26742632.png -- distinguished by a different seed.
|
||||
|
||||
000011.455191342.01.png -- Two files produced by the same command using
|
||||
000011.455191342.02.png -- a batch size>1 (e.g. -b2). They have the same seed.
|
||||
|
||||
000011.4160627868.grid#1-4.png -- a grid of four images (-g); the whole grid can
|
||||
be regenerated with the indicated key
|
||||
|
||||
* It should no longer be possible for one image to overwrite another
|
||||
* You can use the "cd" and "pwd" commands at the dream> prompt to set and retrieve
|
||||
the path of the output directory.
|
||||
|
||||
* v1.04 (22 August 2022 - after the drop)
|
||||
* Updated README to reflect installation of the released weights.
|
||||
* Suppressed very noisy and inconsequential warning when loading the frozen CLIP
|
||||
tokenizer.
|
||||
|
||||
* v1.03 (22 August 2022)
|
||||
* The original txt2img and img2img scripts from the CompViz repository have been moved into
|
||||
a subfolder named "orig_scripts", to reduce confusion.
|
||||
|
||||
* v1.02 (21 August 2022)
|
||||
* A copy of the prompt and all of its switches and options is now stored in the corresponding
|
||||
image in a tEXt metadata field named "Dream". You can read the prompt using scripts/images2prompt.py,
|
||||
or an image editor that allows you to explore the full metadata.
|
||||
**Please run "conda env update -f environment.yaml" to load the k_lms dependencies!!**
|
||||
|
||||
* v1.01 (21 August 2022)
|
||||
* added k_lms sampling.
|
||||
**Please run "conda env update -f environment.yaml" to load the k_lms dependencies!!**
|
||||
* use half precision arithmetic by default, resulting in faster execution and lower memory requirements
|
||||
Pass argument --full_precision to dream.py to get slower but more accurate image generation
|
||||
<div align="center"><img src="docs/assets/invoke-web-server-1.png" width=640></div>
|
||||
|
||||
|
||||
## Installation
|
||||
_Note: This fork is rapidly evolving. Please use the
|
||||
[Issues](https://github.com/invoke-ai/InvokeAI/issues) tab to report bugs and make feature
|
||||
requests. Be sure to use the provided templates. They will help aid diagnose issues faster._
|
||||
|
||||
There are separate installation walkthroughs for [Linux/Mac](#linuxmac) and [Windows](#windows).
|
||||
## Table of Contents
|
||||
|
||||
### Linux/Mac
|
||||
1. [Installation](#installation)
|
||||
2. [Hardware Requirements](#hardware-requirements)
|
||||
3. [Features](#features)
|
||||
4. [Latest Changes](#latest-changes)
|
||||
5. [Troubleshooting](#troubleshooting)
|
||||
6. [Contributing](#contributing)
|
||||
7. [Contributors](#contributors)
|
||||
8. [Support](#support)
|
||||
9. [Further Reading](#further-reading)
|
||||
|
||||
1. You will need to install the following prerequisites if they are not already available. Use your
|
||||
operating system's preferred installer
|
||||
* Python (version 3.8.5 recommended; higher may work)
|
||||
* git
|
||||
### Installation
|
||||
|
||||
2. Install the Python Anaconda environment manager using pip3.
|
||||
```
|
||||
~$ pip3 install anaconda
|
||||
```
|
||||
After installing anaconda, you should log out of your system and log back in. If the installation
|
||||
worked, your command prompt will be prefixed by the name of the current anaconda environment, "(base)".
|
||||
This fork is supported across multiple platforms. You can find individual installation instructions
|
||||
below.
|
||||
|
||||
3. Copy the stable-diffusion source code from GitHub:
|
||||
```
|
||||
(base) ~$ git clone https://github.com/lstein/stable-diffusion.git
|
||||
```
|
||||
This will create stable-diffusion folder where you will follow the rest of the steps.
|
||||
- #### [Linux](docs/installation/INSTALL_LINUX.md)
|
||||
|
||||
4. Enter the newly-created stable-diffusion folder. From this step forward make sure that you are working in the stable-diffusion directory!
|
||||
```
|
||||
(base) ~$ cd stable-diffusion
|
||||
(base) ~/stable-diffusion$
|
||||
```
|
||||
5. Use anaconda to copy necessary python packages, create a new python environment named "ldm",
|
||||
and activate the environment.
|
||||
```
|
||||
(base) ~/stable-diffusion$ conda env create -f environment.yaml
|
||||
(base) ~/stable-diffusion$ conda activate ldm
|
||||
(ldm) ~/stable-diffusion$
|
||||
```
|
||||
After these steps, your command prompt will be prefixed by "(ldm)" as shown above.
|
||||
- #### [Windows](docs/installation/INSTALL_WINDOWS.md)
|
||||
|
||||
6. Load a couple of small machine-learning models required by stable diffusion:
|
||||
```
|
||||
(ldm) ~/stable-diffusion$ python3 scripts/preload_models.py
|
||||
- #### [Macintosh](docs/installation/INSTALL_MAC.md)
|
||||
|
||||
### Hardware Requirements
|
||||
|
||||
#### System
|
||||
|
||||
You wil need one of the following:
|
||||
|
||||
- An NVIDIA-based graphics card with 4 GB or more VRAM memory.
|
||||
- An Apple computer with an M1 chip.
|
||||
|
||||
#### Memory
|
||||
|
||||
- At least 12 GB Main Memory RAM.
|
||||
|
||||
#### Disk
|
||||
|
||||
- At least 12 GB of free disk space for the machine learning model, Python, and all its dependencies.
|
||||
|
||||
**Note**
|
||||
|
||||
If you have a Nvidia 10xx series card (e.g. the 1080ti), please
|
||||
run the dream script in full-precision mode as shown below.
|
||||
|
||||
Similarly, specify full-precision mode on Apple M1 hardware.
|
||||
|
||||
Precision is auto configured based on the device. If however you encounter
|
||||
errors like 'expected type Float but found Half' or 'not implemented for Half'
|
||||
you can try starting `invoke.py` with the `--precision=float32` flag:
|
||||
|
||||
```bash
|
||||
(ldm) ~/stable-diffusion$ python scripts/invoke.py --precision=float32
|
||||
```
|
||||
|
||||
Note that this step is necessary because I modified the original
|
||||
just-in-time model loading scheme to allow the script to work on GPU
|
||||
machines that are not internet connected. See [Workaround for machines with limited internet connectivity](#workaround-for-machines-with-limited-internet-connectivity)
|
||||
### Features
|
||||
|
||||
7. Now you need to install the weights for the stable diffusion model.
|
||||
#### Major Features
|
||||
|
||||
For running with the released weights, you will first need to set up an acount with Hugging Face (https://huggingface.co).
|
||||
Use your credentials to log in, and then point your browser at https://huggingface.co/CompVis/stable-diffusion-v-1-4-original.
|
||||
You may be asked to sign a license agreement at this point.
|
||||
- [Web Server](docs/features/WEB.md)
|
||||
- [Interactive Command Line Interface](docs/features/CLI.md)
|
||||
- [Image To Image](docs/features/IMG2IMG.md)
|
||||
- [Inpainting Support](docs/features/INPAINTING.md)
|
||||
- [Outpainting Support](docs/features/OUTPAINTING.md)
|
||||
- [Upscaling, face-restoration and outpainting](docs/features/POSTPROCESS.md)
|
||||
- [Seamless Tiling](docs/features/OTHER.md#seamless-tiling)
|
||||
- [Google Colab](docs/features/OTHER.md#google-colab)
|
||||
- [Reading Prompts From File](docs/features/PROMPTS.md#reading-prompts-from-a-file)
|
||||
- [Shortcut: Reusing Seeds](docs/features/OTHER.md#shortcuts-reusing-seeds)
|
||||
- [Prompt Blending](docs/features/PROMPTS.md#prompt-blending)
|
||||
- [Thresholding and Perlin Noise Initialization Options](/docs/features/OTHER.md#thresholding-and-perlin-noise-initialization-options)
|
||||
- [Negative/Unconditioned Prompts](docs/features/PROMPTS.md#negative-and-unconditioned-prompts)
|
||||
- [Variations](docs/features/VARIATIONS.md)
|
||||
- [Personalizing Text-to-Image Generation](docs/features/TEXTUAL_INVERSION.md)
|
||||
- [Simplified API for text to image generation](docs/features/OTHER.md#simplified-api)
|
||||
|
||||
Click on "Files and versions" near the top of the page, and then click on the file named "sd-v1-4.ckpt". You'll be taken
|
||||
to a page that prompts you to click the "download" link. Save the file somewhere safe on your local machine.
|
||||
#### Other Features
|
||||
|
||||
Now run the following commands from within the stable-diffusion directory. This will create a symbolic
|
||||
link from the stable-diffusion model.ckpt file, to the true location of the sd-v1-4.ckpt file.
|
||||
|
||||
```
|
||||
(ldm) ~/stable-diffusion$ mkdir -p models/ldm/stable-diffusion-v1
|
||||
(ldm) ~/stable-diffusion$ ln -sf /path/to/sd-v1-4.ckpt models/ldm/stable-diffusion-v1/model.ckpt
|
||||
```
|
||||
- [Creating Transparent Regions for Inpainting](docs/features/INPAINTING.md#creating-transparent-regions-for-inpainting)
|
||||
- [Preload Models](docs/features/OTHER.md#preload-models)
|
||||
|
||||
8. Start generating images!
|
||||
```
|
||||
# for the pre-release weights use the -l or --liaon400m switch
|
||||
(ldm) ~/stable-diffusion$ python3 scripts/dream.py -l
|
||||
### Latest Changes
|
||||
|
||||
# for the post-release weights do not use the switch
|
||||
(ldm) ~/stable-diffusion$ python3 scripts/dream.py
|
||||
- v2.0.1 (13 October 2022)
|
||||
- fix noisy images at high step count when using k* samplers
|
||||
- dream.py script now calls invoke.py module directly rather than
|
||||
via a new python process (which could break the environment)
|
||||
|
||||
# for additional configuration switches and arguments, use -h or --help
|
||||
(ldm) ~/stable-diffusion$ python3 scripts/dream.py -h
|
||||
```
|
||||
9. Subsequently, to relaunch the script, be sure to run "conda activate ldm" (step 5, second command), enter the "stable-diffusion"
|
||||
directory, and then launch the dream script (step 8). If you forget to activate the ldm environment, the script will fail with multiple ModuleNotFound errors.
|
||||
- v2.0.0 (9 October 2022)
|
||||
|
||||
#### Updating to newer versions of the script
|
||||
- `dream.py` script renamed `invoke.py`. A `dream.py` script wrapper remains
|
||||
for backward compatibility.
|
||||
- Completely new WebGUI - launch with `python3 scripts/invoke.py --web`
|
||||
- Support for <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/INPAINTING.md">inpainting</a> and <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/OUTPAINTING.md">outpainting</a>
|
||||
- img2img runs on all k* samplers
|
||||
- Support for <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/PROMPTS.md#negative-and-unconditioned-prompts">negative prompts</a>
|
||||
- Support for CodeFormer face reconstruction
|
||||
- Support for Textual Inversion on Macintoshes
|
||||
- Support in both WebGUI and CLI for <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/POSTPROCESS.md">post-processing of previously-generated images</a>
|
||||
using facial reconstruction, ESRGAN upscaling, outcropping (similar to DALL-E infinite canvas),
|
||||
and "embiggen" upscaling. See the `!fix` command.
|
||||
- New `--hires` option on `invoke>` line allows <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/CLI.md#this-is-an-example-of-txt2img">larger images to be created without duplicating elements</a>, at the cost of some performance.
|
||||
- New `--perlin` and `--threshold` options allow you to add and control variation
|
||||
during image generation (see <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/OTHER.md#thresholding-and-perlin-noise-initialization-options">Thresholding and Perlin Noise Initialization</a>
|
||||
- Extensive metadata now written into PNG files, allowing reliable regeneration of images
|
||||
and tweaking of previous settings.
|
||||
- Command-line completion in `invoke.py` now works on Windows, Linux and Mac platforms.
|
||||
- Improved <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/CLI.md">command-line completion behavior</a>.
|
||||
New commands added:
|
||||
* List command-line history with `!history`
|
||||
* Search command-line history with `!search`
|
||||
* Clear history with `!clear`
|
||||
- Deprecated `--full_precision` / `-F`. Simply omit it and `invoke.py` will auto
|
||||
configure. To switch away from auto use the new flag like `--precision=float32`.
|
||||
|
||||
This distribution is changing rapidly. If you used the "git clone" method (step 5) to download the stable-diffusion directory, then to update to the latest and greatest version, launch the Anaconda window, enter "stable-diffusion", and type:
|
||||
```
|
||||
(ldm) ~/stable-diffusion$ git pull
|
||||
```
|
||||
This will bring your local copy into sync with the remote one.
|
||||
For older changelogs, please visit the **[CHANGELOG](docs/features/CHANGELOG.md)**.
|
||||
|
||||
### Windows
|
||||
### Troubleshooting
|
||||
|
||||
1. Install Anaconda3 (miniconda3 version) from here: https://docs.anaconda.com/anaconda/install/windows/
|
||||
Please check out our **[Q&A](docs/help/TROUBLESHOOT.md)** to get solutions for common installation
|
||||
problems and other issues.
|
||||
|
||||
2. Install Git from here: https://git-scm.com/download/win
|
||||
# Contributing
|
||||
|
||||
3. Launch Anaconda from the Windows Start menu. This will bring up a command window. Type all the remaining commands in this window.
|
||||
Anyone who wishes to contribute to this project, whether documentation, features, bug fixes, code
|
||||
cleanup, testing, or code reviews, is very much encouraged to do so. If you are unfamiliar with how
|
||||
to contribute to GitHub projects, here is a
|
||||
[Getting Started Guide](https://opensource.com/article/19/7/create-pull-request-github).
|
||||
|
||||
4. Run the command:
|
||||
```
|
||||
git clone https://github.com/lstein/stable-diffusion.git
|
||||
```
|
||||
This will create stable-diffusion folder where you will follow the rest of the steps.
|
||||
A full set of contribution guidelines, along with templates, are in progress, but for now the most
|
||||
important thing is to **make your pull request against the "development" branch**, and not against
|
||||
"main". This will help keep public breakage to a minimum and will allow you to propose more radical
|
||||
changes.
|
||||
|
||||
5. Enter the newly-created stable-diffusion folder. From this step forward make sure that you are working in the stable-diffusion directory!
|
||||
```
|
||||
cd stable-diffusion
|
||||
```
|
||||
### Contributors
|
||||
|
||||
6. Run the following two commands:
|
||||
```
|
||||
conda env create -f environment.yaml (step 6a)
|
||||
conda activate ldm (step 6b)
|
||||
```
|
||||
This will install all python requirements and activate the "ldm" environment which sets PATH and other environment variables properly.
|
||||
This fork is a combined effort of various people from across the world.
|
||||
[Check out the list of all these amazing people](docs/other/CONTRIBUTORS.md). We thank them for
|
||||
their time, hard work and effort.
|
||||
|
||||
7. Run the command:
|
||||
```
|
||||
python scripts\preload_models.py
|
||||
```
|
||||
### Support
|
||||
|
||||
This installs several machine learning models that stable diffusion
|
||||
requires. (Note that this step is required. I created it because some people
|
||||
are using GPU systems that are behind a firewall and the models can't be
|
||||
downloaded just-in-time)
|
||||
For support, please use this repository's GitHub Issues tracking service. Feel free to send me an
|
||||
email if you use and like the script.
|
||||
|
||||
8. Now you need to install the weights for the big stable diffusion model.
|
||||
Original portions of the software are Copyright (c) 2020
|
||||
[Lincoln D. Stein](https://github.com/lstein)
|
||||
|
||||
For running with the released weights, you will first need to set up
|
||||
an acount with Hugging Face (https://huggingface.co). Use your
|
||||
credentials to log in, and then point your browser at
|
||||
https://huggingface.co/CompVis/stable-diffusion-v-1-4-original. You
|
||||
may be asked to sign a license agreement at this point.
|
||||
### Further Reading
|
||||
|
||||
Click on "Files and versions" near the top of the page, and then click
|
||||
on the file named "sd-v1-4.ckpt". You'll be taken to a page that
|
||||
prompts you to click the "download" link. Now save the file somewhere
|
||||
safe on your local machine. The weight file is >4 GB in size, so
|
||||
downloading may take a while.
|
||||
|
||||
Now run the following commands from **within the stable-diffusion
|
||||
directory** to copy the weights file to the right place:
|
||||
|
||||
```
|
||||
mkdir -p models\ldm\stable-diffusion-v1
|
||||
copy C:\path\to\sd-v1-4.ckpt models\ldm\stable-diffusion-v1\model.ckpt
|
||||
```
|
||||
Please replace "C:\path\to\sd-v1.4.ckpt" with the correct path to wherever
|
||||
you stashed this file. If you prefer not to copy or move the .ckpt file,
|
||||
you may instead create a shortcut to it from within
|
||||
"models\ldm\stable-diffusion-v1\".
|
||||
|
||||
9. Start generating images!
|
||||
```
|
||||
# for the pre-release weights
|
||||
python scripts\dream.py -l
|
||||
|
||||
# for the post-release weights
|
||||
python scripts\dream.py
|
||||
```
|
||||
10. Subsequently, to relaunch the script, first activate the Anaconda command window (step 3), enter the stable-diffusion directory (step 5, "cd \path\to\stable-diffusion"), run "conda activate ldm" (step 6b), and then launch the dream script (step 9).
|
||||
|
||||
#### Updating to newer versions of the script
|
||||
|
||||
This distribution is changing rapidly. If you used the "git clone" method (step 5) to download the stable-diffusion directory, then to update to the latest and greatest version, launch the Anaconda window, enter "stable-diffusion", and type:
|
||||
```
|
||||
git pull
|
||||
```
|
||||
This will bring your local copy into sync with the remote one.
|
||||
|
||||
## Simplified API for text to image generation
|
||||
|
||||
For programmers who wish to incorporate stable-diffusion into other
|
||||
products, this repository includes a simplified API for text to image generation, which
|
||||
lets you create images from a prompt in just three lines of code:
|
||||
|
||||
~~~~
|
||||
from ldm.simplet2i import T2I
|
||||
model = T2I()
|
||||
outputs = model.txt2img("a unicorn in manhattan")
|
||||
~~~~
|
||||
|
||||
Outputs is a list of lists in the format [[filename1,seed1],[filename2,seed2]...]
|
||||
Please see ldm/simplet2i.py for more information.
|
||||
|
||||
|
||||
## Workaround for machines with limited internet connectivity
|
||||
|
||||
My development machine is a GPU node in a high-performance compute
|
||||
cluster which has no connection to the internet. During model
|
||||
initialization, stable-diffusion tries to download the Bert tokenizer
|
||||
and a file needed by the kornia library. This obviously didn't work
|
||||
for me.
|
||||
|
||||
To work around this, I have modified ldm/modules/encoders/modules.py
|
||||
to look for locally cached Bert files rather than attempting to
|
||||
download them. For this to work, you must run
|
||||
"scripts/preload_models.py" once from an internet-connected machine
|
||||
prior to running the code on an isolated one. This assumes that both
|
||||
machines share a common network-mounted filesystem with a common
|
||||
.cache directory.
|
||||
|
||||
~~~~
|
||||
(ldm) ~/stable-diffusion$ python3 ./scripts/preload_models.py
|
||||
preloading bert tokenizer...
|
||||
Downloading: 100%|██████████████████████████████████| 28.0/28.0 [00:00<00:00, 49.3kB/s]
|
||||
Downloading: 100%|██████████████████████████████████| 226k/226k [00:00<00:00, 2.79MB/s]
|
||||
Downloading: 100%|██████████████████████████████████| 455k/455k [00:00<00:00, 4.36MB/s]
|
||||
Downloading: 100%|██████████████████████████████████| 570/570 [00:00<00:00, 477kB/s]
|
||||
...success
|
||||
preloading kornia requirements...
|
||||
Downloading: "https://github.com/DagnyT/hardnet/raw/master/pretrained/train_liberty_with_aug/checkpoint_liberty_with_aug.pth" to /u/lstein/.cache/torch/hub/checkpoints/checkpoint_liberty_with_aug.pth
|
||||
100%|███████████████████████████████████████████████| 5.10M/5.10M [00:00<00:00, 101MB/s]
|
||||
...success
|
||||
~~~~
|
||||
|
||||
If you don't need this change and want to download the files just in
|
||||
time, copy over the file ldm/modules/encoders/modules.py from the
|
||||
CompVis/stable-diffusion repository. Or you can run preload_models.py
|
||||
on the target machine.
|
||||
|
||||
## Support
|
||||
|
||||
For support,
|
||||
please use this repository's GitHub Issues tracking service. Feel free
|
||||
to send me an email if you use and like the script.
|
||||
|
||||
*Original Author:* Lincoln D. Stein <lincoln.stein@gmail.com>
|
||||
|
||||
*Contributions by:*
|
||||
[Peter Kowalczyk](https://github.com/slix), [Henry Harrison](https://github.com/hwharrison),
|
||||
[xraxra](https://github.com/xraxra), [bmaltais](https://github.com/bmaltais), [Sean McLellan](https://github.com/Oceanswave),
|
||||
[nicolai256](https://github.com/nicolai256), [Benjamin Warner](https://github.com/warner-benjamin),
|
||||
[tildebyte](https://github.com/tildebyte),
|
||||
and [Tesseract Cat](https://github.com/TesseractCat)
|
||||
|
||||
|
||||
Original portions of the software are Copyright (c) 2020 Lincoln D. Stein (https://github.com/lstein)
|
||||
|
||||
#Further Reading
|
||||
|
||||
Please see the original README for more information on this software
|
||||
and underlying algorithm, located in the file README-CompViz.md.
|
||||
Please see the original README for more information on this software and underlying algorithm,
|
||||
located in the file [README-CompViz.md](docs/other/README-CompViz.md).
|
||||
|
||||
31
TODO.txt
@@ -1,31 +0,0 @@
|
||||
Feature requests:
|
||||
|
||||
1. "gobig" mode - split image into strips, scale up, add detail using
|
||||
img2img and reassemble with feathering. Issue #66.
|
||||
|
||||
2. Port basujindal low VRAM optimizations. Issue #62
|
||||
|
||||
3. Store images under folders named after the prompt. Issue #27.
|
||||
|
||||
4. Some sort of automation for generating variations. Issues #32 and #47.
|
||||
|
||||
5. Support for inpainting masks #68.
|
||||
|
||||
6. Support for loading variations of the stable-diffusion
|
||||
weights #49
|
||||
|
||||
7. Support for klms and other non-ddim samplers in img2img() #36
|
||||
|
||||
8. Pass a shell command to open up an image viewer on the last
|
||||
batch of images generated #29.
|
||||
|
||||
Code Refactorization:
|
||||
|
||||
1. Move the PNG file generation code out of simplet2i and into
|
||||
separate module. txt2img() and img2img() should return Image
|
||||
objects, and parent code is responsible for filenaming logic.
|
||||
|
||||
2. Refactor redundant code that is shared between txt2img() and
|
||||
img2img().
|
||||
|
||||
3. Experiment with replacing CompViz code with HuggingFace.
|
||||
1083
backend/invoke_ai_web_server.py
Normal file
55
backend/modules/create_cmd_parser.py
Normal file
@@ -0,0 +1,55 @@
|
||||
import argparse
|
||||
import os
|
||||
from ldm.invoke.args import PRECISION_CHOICES
|
||||
|
||||
|
||||
def create_cmd_parser():
|
||||
parser = argparse.ArgumentParser(description="InvokeAI web UI")
|
||||
parser.add_argument(
|
||||
"--host",
|
||||
type=str,
|
||||
help="The host to serve on",
|
||||
default="localhost",
|
||||
)
|
||||
parser.add_argument("--port", type=int, help="The port to serve on", default=9090)
|
||||
parser.add_argument(
|
||||
"--cors",
|
||||
nargs="*",
|
||||
type=str,
|
||||
help="Additional allowed origins, comma-separated",
|
||||
)
|
||||
parser.add_argument(
|
||||
"--embedding_path",
|
||||
type=str,
|
||||
help="Path to a pre-trained embedding manager checkpoint - can only be set on command line",
|
||||
)
|
||||
# TODO: Can't get flask to serve images from any dir (saving to the dir does work when specified)
|
||||
# parser.add_argument(
|
||||
# "--output_dir",
|
||||
# default="outputs/",
|
||||
# type=str,
|
||||
# help="Directory for output images",
|
||||
# )
|
||||
parser.add_argument(
|
||||
"-v",
|
||||
"--verbose",
|
||||
action="store_true",
|
||||
help="Enables verbose logging",
|
||||
)
|
||||
parser.add_argument(
|
||||
"--precision",
|
||||
dest="precision",
|
||||
type=str,
|
||||
choices=PRECISION_CHOICES,
|
||||
metavar="PRECISION",
|
||||
help=f'Set model precision. Defaults to auto selected based on device. Options: {", ".join(PRECISION_CHOICES)}',
|
||||
default="auto",
|
||||
)
|
||||
parser.add_argument(
|
||||
'--free_gpu_mem',
|
||||
dest='free_gpu_mem',
|
||||
action='store_true',
|
||||
help='Force free gpu memory before final decoding',
|
||||
)
|
||||
|
||||
return parser
|
||||
61
backend/modules/parameters.py
Normal file
@@ -0,0 +1,61 @@
|
||||
from backend.modules.parse_seed_weights import parse_seed_weights
|
||||
import argparse
|
||||
|
||||
SAMPLER_CHOICES = [
|
||||
"ddim",
|
||||
"k_dpm_2_a",
|
||||
"k_dpm_2",
|
||||
"k_euler_a",
|
||||
"k_euler",
|
||||
"k_heun",
|
||||
"k_lms",
|
||||
"plms",
|
||||
]
|
||||
|
||||
|
||||
def parameters_to_command(params):
|
||||
"""
|
||||
Converts dict of parameters into a `invoke.py` REPL command.
|
||||
"""
|
||||
|
||||
switches = list()
|
||||
|
||||
if "prompt" in params:
|
||||
switches.append(f'"{params["prompt"]}"')
|
||||
if "steps" in params:
|
||||
switches.append(f'-s {params["steps"]}')
|
||||
if "seed" in params:
|
||||
switches.append(f'-S {params["seed"]}')
|
||||
if "width" in params:
|
||||
switches.append(f'-W {params["width"]}')
|
||||
if "height" in params:
|
||||
switches.append(f'-H {params["height"]}')
|
||||
if "cfg_scale" in params:
|
||||
switches.append(f'-C {params["cfg_scale"]}')
|
||||
if "sampler_name" in params:
|
||||
switches.append(f'-A {params["sampler_name"]}')
|
||||
if "seamless" in params and params["seamless"] == True:
|
||||
switches.append(f"--seamless")
|
||||
if "init_img" in params and len(params["init_img"]) > 0:
|
||||
switches.append(f'-I {params["init_img"]}')
|
||||
if "init_mask" in params and len(params["init_mask"]) > 0:
|
||||
switches.append(f'-M {params["init_mask"]}')
|
||||
if "init_color" in params and len(params["init_color"]) > 0:
|
||||
switches.append(f'--init_color {params["init_color"]}')
|
||||
if "strength" in params and "init_img" in params:
|
||||
switches.append(f'-f {params["strength"]}')
|
||||
if "fit" in params and params["fit"] == True:
|
||||
switches.append(f"--fit")
|
||||
if "gfpgan_strength" in params and params["gfpgan_strength"]:
|
||||
switches.append(f'-G {params["gfpgan_strength"]}')
|
||||
if "upscale" in params and params["upscale"]:
|
||||
switches.append(f'-U {params["upscale"][0]} {params["upscale"][1]}')
|
||||
if "variation_amount" in params and params["variation_amount"] > 0:
|
||||
switches.append(f'-v {params["variation_amount"]}')
|
||||
if "with_variations" in params:
|
||||
seed_weight_pairs = ",".join(
|
||||
f"{seed}:{weight}" for seed, weight in params["with_variations"]
|
||||
)
|
||||
switches.append(f"-V {seed_weight_pairs}")
|
||||
|
||||
return " ".join(switches)
|
||||
47
backend/modules/parse_seed_weights.py
Normal file
@@ -0,0 +1,47 @@
|
||||
def parse_seed_weights(seed_weights):
|
||||
"""
|
||||
Accepts seed weights as string in "12345:0.1,23456:0.2,3456:0.3" format
|
||||
Validates them
|
||||
If valid: returns as [[12345, 0.1], [23456, 0.2], [3456, 0.3]]
|
||||
If invalid: returns False
|
||||
"""
|
||||
|
||||
# Must be a string
|
||||
if not isinstance(seed_weights, str):
|
||||
return False
|
||||
# String must not be empty
|
||||
if len(seed_weights) == 0:
|
||||
return False
|
||||
|
||||
pairs = []
|
||||
|
||||
for pair in seed_weights.split(","):
|
||||
split_values = pair.split(":")
|
||||
|
||||
# Seed and weight are required
|
||||
if len(split_values) != 2:
|
||||
return False
|
||||
|
||||
if len(split_values[0]) == 0 or len(split_values[1]) == 1:
|
||||
return False
|
||||
|
||||
# Try casting the seed to int and weight to float
|
||||
try:
|
||||
seed = int(split_values[0])
|
||||
weight = float(split_values[1])
|
||||
except ValueError:
|
||||
return False
|
||||
|
||||
# Seed must be 0 or above
|
||||
if not seed >= 0:
|
||||
return False
|
||||
|
||||
# Weight must be between 0 and 1
|
||||
if not (weight >= 0 and weight <= 1):
|
||||
return False
|
||||
|
||||
# This pair is valid
|
||||
pairs.append([seed, weight])
|
||||
|
||||
# All pairs are valid
|
||||
return pairs
|
||||
821
backend/server.py
Normal file
@@ -0,0 +1,821 @@
|
||||
import mimetypes
|
||||
import transformers
|
||||
import json
|
||||
import os
|
||||
import traceback
|
||||
import eventlet
|
||||
import glob
|
||||
import shlex
|
||||
import math
|
||||
import shutil
|
||||
import sys
|
||||
|
||||
sys.path.append(".")
|
||||
|
||||
from argparse import ArgumentTypeError
|
||||
from modules.create_cmd_parser import create_cmd_parser
|
||||
|
||||
parser = create_cmd_parser()
|
||||
opt = parser.parse_args()
|
||||
|
||||
|
||||
from flask_socketio import SocketIO
|
||||
from flask import Flask, send_from_directory, url_for, jsonify
|
||||
from pathlib import Path
|
||||
from PIL import Image
|
||||
from pytorch_lightning import logging
|
||||
from threading import Event
|
||||
from uuid import uuid4
|
||||
from send2trash import send2trash
|
||||
|
||||
|
||||
from ldm.generate import Generate
|
||||
from ldm.invoke.restoration import Restoration
|
||||
from ldm.invoke.pngwriter import PngWriter, retrieve_metadata
|
||||
from ldm.invoke.args import APP_ID, APP_VERSION, calculate_init_img_hash
|
||||
from ldm.invoke.conditioning import split_weighted_subprompts
|
||||
|
||||
from modules.parameters import parameters_to_command
|
||||
|
||||
|
||||
"""
|
||||
USER CONFIG
|
||||
"""
|
||||
if opt.cors and "*" in opt.cors:
|
||||
raise ArgumentTypeError('"*" is not an allowed CORS origin')
|
||||
|
||||
|
||||
output_dir = "outputs/" # Base output directory for images
|
||||
host = opt.host # Web & socket.io host
|
||||
port = opt.port # Web & socket.io port
|
||||
verbose = opt.verbose # enables copious socket.io logging
|
||||
precision = opt.precision
|
||||
free_gpu_mem = opt.free_gpu_mem
|
||||
embedding_path = opt.embedding_path
|
||||
additional_allowed_origins = (
|
||||
opt.cors if opt.cors else []
|
||||
) # additional CORS allowed origins
|
||||
model = "stable-diffusion-1.4"
|
||||
|
||||
"""
|
||||
END USER CONFIG
|
||||
"""
|
||||
|
||||
|
||||
print("* Initializing, be patient...\n")
|
||||
|
||||
|
||||
"""
|
||||
SERVER SETUP
|
||||
"""
|
||||
|
||||
|
||||
# fix missing mimetypes on windows due to registry wonkiness
|
||||
mimetypes.add_type("application/javascript", ".js")
|
||||
mimetypes.add_type("text/css", ".css")
|
||||
|
||||
app = Flask(__name__, static_url_path="", static_folder="../frontend/dist/")
|
||||
|
||||
|
||||
app.config["OUTPUTS_FOLDER"] = "../outputs"
|
||||
|
||||
|
||||
@app.route("/outputs/<path:filename>")
|
||||
def outputs(filename):
|
||||
return send_from_directory(app.config["OUTPUTS_FOLDER"], filename)
|
||||
|
||||
|
||||
@app.route("/", defaults={"path": ""})
|
||||
def serve(path):
|
||||
return send_from_directory(app.static_folder, "index.html")
|
||||
|
||||
|
||||
logger = True if verbose else False
|
||||
engineio_logger = True if verbose else False
|
||||
|
||||
# default 1,000,000, needs to be higher for socketio to accept larger images
|
||||
max_http_buffer_size = 10000000
|
||||
|
||||
cors_allowed_origins = [f"http://{host}:{port}"] + additional_allowed_origins
|
||||
|
||||
socketio = SocketIO(
|
||||
app,
|
||||
logger=logger,
|
||||
engineio_logger=engineio_logger,
|
||||
max_http_buffer_size=max_http_buffer_size,
|
||||
cors_allowed_origins=cors_allowed_origins,
|
||||
ping_interval=(50, 50),
|
||||
ping_timeout=60,
|
||||
)
|
||||
|
||||
|
||||
"""
|
||||
END SERVER SETUP
|
||||
"""
|
||||
|
||||
|
||||
"""
|
||||
APP SETUP
|
||||
"""
|
||||
|
||||
|
||||
class CanceledException(Exception):
|
||||
pass
|
||||
|
||||
|
||||
try:
|
||||
gfpgan, codeformer, esrgan = None, None, None
|
||||
from ldm.invoke.restoration.base import Restoration
|
||||
|
||||
restoration = Restoration()
|
||||
gfpgan, codeformer = restoration.load_face_restore_models()
|
||||
esrgan = restoration.load_esrgan()
|
||||
|
||||
# coreformer.process(self, image, strength, device, seed=None, fidelity=0.75)
|
||||
|
||||
except (ModuleNotFoundError, ImportError):
|
||||
print(traceback.format_exc(), file=sys.stderr)
|
||||
print(">> You may need to install the ESRGAN and/or GFPGAN modules")
|
||||
|
||||
canceled = Event()
|
||||
|
||||
# reduce logging outputs to error
|
||||
transformers.logging.set_verbosity_error()
|
||||
logging.getLogger("pytorch_lightning").setLevel(logging.ERROR)
|
||||
|
||||
# Initialize and load model
|
||||
generate = Generate(
|
||||
model,
|
||||
precision=precision,
|
||||
embedding_path=embedding_path,
|
||||
)
|
||||
generate.free_gpu_mem = free_gpu_mem
|
||||
generate.load_model()
|
||||
|
||||
|
||||
# location for "finished" images
|
||||
result_path = os.path.join(output_dir, "img-samples/")
|
||||
|
||||
# temporary path for intermediates
|
||||
intermediate_path = os.path.join(result_path, "intermediates/")
|
||||
|
||||
# path for user-uploaded init images and masks
|
||||
init_image_path = os.path.join(result_path, "init-images/")
|
||||
mask_image_path = os.path.join(result_path, "mask-images/")
|
||||
|
||||
# txt log
|
||||
log_path = os.path.join(result_path, "invoke_log.txt")
|
||||
|
||||
# make all output paths
|
||||
[
|
||||
os.makedirs(path, exist_ok=True)
|
||||
for path in [result_path, intermediate_path, init_image_path, mask_image_path]
|
||||
]
|
||||
|
||||
|
||||
"""
|
||||
END APP SETUP
|
||||
"""
|
||||
|
||||
|
||||
"""
|
||||
SOCKET.IO LISTENERS
|
||||
"""
|
||||
|
||||
|
||||
@socketio.on("requestSystemConfig")
|
||||
def handle_request_capabilities():
|
||||
print(f">> System config requested")
|
||||
config = get_system_config()
|
||||
socketio.emit("systemConfig", config)
|
||||
|
||||
|
||||
@socketio.on("requestImages")
|
||||
def handle_request_images(page=1, offset=0, last_mtime=None):
|
||||
chunk_size = 50
|
||||
|
||||
if last_mtime:
|
||||
print(f">> Latest images requested")
|
||||
else:
|
||||
print(
|
||||
f">> Page {page} of images requested (page size {chunk_size} offset {offset})"
|
||||
)
|
||||
|
||||
paths = glob.glob(os.path.join(result_path, "*.png"))
|
||||
sorted_paths = sorted(paths, key=lambda x: os.path.getmtime(x), reverse=True)
|
||||
|
||||
if last_mtime:
|
||||
image_paths = filter(lambda x: os.path.getmtime(x) > last_mtime, sorted_paths)
|
||||
else:
|
||||
|
||||
image_paths = sorted_paths[
|
||||
slice(chunk_size * (page - 1) + offset, chunk_size * page + offset)
|
||||
]
|
||||
page = page + 1
|
||||
|
||||
image_array = []
|
||||
|
||||
for path in image_paths:
|
||||
metadata = retrieve_metadata(path)
|
||||
image_array.append(
|
||||
{
|
||||
"url": path,
|
||||
"mtime": os.path.getmtime(path),
|
||||
"metadata": metadata["sd-metadata"],
|
||||
}
|
||||
)
|
||||
|
||||
socketio.emit(
|
||||
"galleryImages",
|
||||
{
|
||||
"images": image_array,
|
||||
"nextPage": page,
|
||||
"offset": offset,
|
||||
"onlyNewImages": True if last_mtime else False,
|
||||
},
|
||||
)
|
||||
|
||||
|
||||
@socketio.on("generateImage")
|
||||
def handle_generate_image_event(
|
||||
generation_parameters, esrgan_parameters, gfpgan_parameters
|
||||
):
|
||||
print(
|
||||
f">> Image generation requested: {generation_parameters}\nESRGAN parameters: {esrgan_parameters}\nGFPGAN parameters: {gfpgan_parameters}"
|
||||
)
|
||||
generate_images(generation_parameters, esrgan_parameters, gfpgan_parameters)
|
||||
|
||||
|
||||
@socketio.on("runESRGAN")
|
||||
def handle_run_esrgan_event(original_image, esrgan_parameters):
|
||||
print(
|
||||
f'>> ESRGAN upscale requested for "{original_image["url"]}": {esrgan_parameters}'
|
||||
)
|
||||
progress = {
|
||||
"currentStep": 1,
|
||||
"totalSteps": 1,
|
||||
"currentIteration": 1,
|
||||
"totalIterations": 1,
|
||||
"currentStatus": "Preparing",
|
||||
"isProcessing": True,
|
||||
"currentStatusHasSteps": False,
|
||||
}
|
||||
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
image = Image.open(original_image["url"])
|
||||
|
||||
seed = (
|
||||
original_image["metadata"]["seed"]
|
||||
if "seed" in original_image["metadata"]
|
||||
else "unknown_seed"
|
||||
)
|
||||
|
||||
progress["currentStatus"] = "Upscaling"
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
image = esrgan.process(
|
||||
image=image,
|
||||
upsampler_scale=esrgan_parameters["upscale"][0],
|
||||
strength=esrgan_parameters["upscale"][1],
|
||||
seed=seed,
|
||||
)
|
||||
|
||||
progress["currentStatus"] = "Saving image"
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
esrgan_parameters["seed"] = seed
|
||||
metadata = parameters_to_post_processed_image_metadata(
|
||||
parameters=esrgan_parameters,
|
||||
original_image_path=original_image["url"],
|
||||
type="esrgan",
|
||||
)
|
||||
command = parameters_to_command(esrgan_parameters)
|
||||
|
||||
path = save_image(image, command, metadata, result_path, postprocessing="esrgan")
|
||||
|
||||
write_log_message(f'[Upscaled] "{original_image["url"]}" > "{path}": {command}')
|
||||
|
||||
progress["currentStatus"] = "Finished"
|
||||
progress["currentStep"] = 0
|
||||
progress["totalSteps"] = 0
|
||||
progress["currentIteration"] = 0
|
||||
progress["totalIterations"] = 0
|
||||
progress["isProcessing"] = False
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
socketio.emit(
|
||||
"esrganResult",
|
||||
{
|
||||
"url": os.path.relpath(path),
|
||||
"mtime": os.path.getmtime(path),
|
||||
"metadata": metadata,
|
||||
},
|
||||
)
|
||||
|
||||
|
||||
@socketio.on("runGFPGAN")
|
||||
def handle_run_gfpgan_event(original_image, gfpgan_parameters):
|
||||
print(
|
||||
f'>> GFPGAN face fix requested for "{original_image["url"]}": {gfpgan_parameters}'
|
||||
)
|
||||
progress = {
|
||||
"currentStep": 1,
|
||||
"totalSteps": 1,
|
||||
"currentIteration": 1,
|
||||
"totalIterations": 1,
|
||||
"currentStatus": "Preparing",
|
||||
"isProcessing": True,
|
||||
"currentStatusHasSteps": False,
|
||||
}
|
||||
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
image = Image.open(original_image["url"])
|
||||
|
||||
seed = (
|
||||
original_image["metadata"]["seed"]
|
||||
if "seed" in original_image["metadata"]
|
||||
else "unknown_seed"
|
||||
)
|
||||
|
||||
progress["currentStatus"] = "Fixing faces"
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
image = gfpgan.process(
|
||||
image=image, strength=gfpgan_parameters["gfpgan_strength"], seed=seed
|
||||
)
|
||||
|
||||
progress["currentStatus"] = "Saving image"
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
gfpgan_parameters["seed"] = seed
|
||||
metadata = parameters_to_post_processed_image_metadata(
|
||||
parameters=gfpgan_parameters,
|
||||
original_image_path=original_image["url"],
|
||||
type="gfpgan",
|
||||
)
|
||||
command = parameters_to_command(gfpgan_parameters)
|
||||
|
||||
path = save_image(image, command, metadata, result_path, postprocessing="gfpgan")
|
||||
|
||||
write_log_message(f'[Fixed faces] "{original_image["url"]}" > "{path}": {command}')
|
||||
|
||||
progress["currentStatus"] = "Finished"
|
||||
progress["currentStep"] = 0
|
||||
progress["totalSteps"] = 0
|
||||
progress["currentIteration"] = 0
|
||||
progress["totalIterations"] = 0
|
||||
progress["isProcessing"] = False
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
socketio.emit(
|
||||
"gfpganResult",
|
||||
{
|
||||
"url": os.path.relpath(path),
|
||||
"mtime": os.path.mtime(path),
|
||||
"metadata": metadata,
|
||||
},
|
||||
)
|
||||
|
||||
|
||||
@socketio.on("cancel")
|
||||
def handle_cancel():
|
||||
print(f">> Cancel processing requested")
|
||||
canceled.set()
|
||||
socketio.emit("processingCanceled")
|
||||
|
||||
|
||||
# TODO: I think this needs a safety mechanism.
|
||||
@socketio.on("deleteImage")
|
||||
def handle_delete_image(path, uuid):
|
||||
print(f'>> Delete requested "{path}"')
|
||||
send2trash(path)
|
||||
socketio.emit("imageDeleted", {"url": path, "uuid": uuid})
|
||||
|
||||
|
||||
# TODO: I think this needs a safety mechanism.
|
||||
@socketio.on("uploadInitialImage")
|
||||
def handle_upload_initial_image(bytes, name):
|
||||
print(f'>> Init image upload requested "{name}"')
|
||||
uuid = uuid4().hex
|
||||
split = os.path.splitext(name)
|
||||
name = f"{split[0]}.{uuid}{split[1]}"
|
||||
file_path = os.path.join(init_image_path, name)
|
||||
os.makedirs(os.path.dirname(file_path), exist_ok=True)
|
||||
newFile = open(file_path, "wb")
|
||||
newFile.write(bytes)
|
||||
socketio.emit("initialImageUploaded", {"url": file_path, "uuid": ""})
|
||||
|
||||
|
||||
# TODO: I think this needs a safety mechanism.
|
||||
@socketio.on("uploadMaskImage")
|
||||
def handle_upload_mask_image(bytes, name):
|
||||
print(f'>> Mask image upload requested "{name}"')
|
||||
uuid = uuid4().hex
|
||||
split = os.path.splitext(name)
|
||||
name = f"{split[0]}.{uuid}{split[1]}"
|
||||
file_path = os.path.join(mask_image_path, name)
|
||||
os.makedirs(os.path.dirname(file_path), exist_ok=True)
|
||||
newFile = open(file_path, "wb")
|
||||
newFile.write(bytes)
|
||||
socketio.emit("maskImageUploaded", {"url": file_path, "uuid": ""})
|
||||
|
||||
|
||||
"""
|
||||
END SOCKET.IO LISTENERS
|
||||
"""
|
||||
|
||||
|
||||
"""
|
||||
ADDITIONAL FUNCTIONS
|
||||
"""
|
||||
|
||||
|
||||
def get_system_config():
|
||||
return {
|
||||
"model": "stable diffusion",
|
||||
"model_id": model,
|
||||
"model_hash": generate.model_hash,
|
||||
"app_id": APP_ID,
|
||||
"app_version": APP_VERSION,
|
||||
}
|
||||
|
||||
|
||||
def parameters_to_post_processed_image_metadata(parameters, original_image_path, type):
|
||||
# top-level metadata minus `image` or `images`
|
||||
metadata = get_system_config()
|
||||
|
||||
orig_hash = calculate_init_img_hash(original_image_path)
|
||||
|
||||
image = {"orig_path": original_image_path, "orig_hash": orig_hash}
|
||||
|
||||
if type == "esrgan":
|
||||
image["type"] = "esrgan"
|
||||
image["scale"] = parameters["upscale"][0]
|
||||
image["strength"] = parameters["upscale"][1]
|
||||
elif type == "gfpgan":
|
||||
image["type"] = "gfpgan"
|
||||
image["strength"] = parameters["gfpgan_strength"]
|
||||
else:
|
||||
raise TypeError(f"Invalid type: {type}")
|
||||
|
||||
metadata["image"] = image
|
||||
return metadata
|
||||
|
||||
|
||||
def parameters_to_generated_image_metadata(parameters):
|
||||
# top-level metadata minus `image` or `images`
|
||||
|
||||
metadata = get_system_config()
|
||||
# remove any image keys not mentioned in RFC #266
|
||||
rfc266_img_fields = [
|
||||
"type",
|
||||
"postprocessing",
|
||||
"sampler",
|
||||
"prompt",
|
||||
"seed",
|
||||
"variations",
|
||||
"steps",
|
||||
"cfg_scale",
|
||||
"threshold",
|
||||
"perlin",
|
||||
"step_number",
|
||||
"width",
|
||||
"height",
|
||||
"extra",
|
||||
"seamless",
|
||||
]
|
||||
|
||||
rfc_dict = {}
|
||||
|
||||
for item in parameters.items():
|
||||
key, value = item
|
||||
if key in rfc266_img_fields:
|
||||
rfc_dict[key] = value
|
||||
|
||||
postprocessing = []
|
||||
|
||||
# 'postprocessing' is either null or an
|
||||
if "gfpgan_strength" in parameters:
|
||||
|
||||
postprocessing.append(
|
||||
{"type": "gfpgan", "strength": float(parameters["gfpgan_strength"])}
|
||||
)
|
||||
|
||||
if "upscale" in parameters:
|
||||
postprocessing.append(
|
||||
{
|
||||
"type": "esrgan",
|
||||
"scale": int(parameters["upscale"][0]),
|
||||
"strength": float(parameters["upscale"][1]),
|
||||
}
|
||||
)
|
||||
|
||||
rfc_dict["postprocessing"] = postprocessing if len(postprocessing) > 0 else None
|
||||
|
||||
# semantic drift
|
||||
rfc_dict["sampler"] = parameters["sampler_name"]
|
||||
|
||||
# display weighted subprompts (liable to change)
|
||||
subprompts = split_weighted_subprompts(parameters["prompt"])
|
||||
subprompts = [{"prompt": x[0], "weight": x[1]} for x in subprompts]
|
||||
rfc_dict["prompt"] = subprompts
|
||||
|
||||
# 'variations' should always exist and be an array, empty or consisting of {'seed': seed, 'weight': weight} pairs
|
||||
variations = []
|
||||
|
||||
if "with_variations" in parameters:
|
||||
variations = [
|
||||
{"seed": x[0], "weight": x[1]} for x in parameters["with_variations"]
|
||||
]
|
||||
|
||||
rfc_dict["variations"] = variations
|
||||
|
||||
if "init_img" in parameters:
|
||||
rfc_dict["type"] = "img2img"
|
||||
rfc_dict["strength"] = parameters["strength"]
|
||||
rfc_dict["fit"] = parameters["fit"] # TODO: Noncompliant
|
||||
rfc_dict["orig_hash"] = calculate_init_img_hash(parameters["init_img"])
|
||||
rfc_dict["init_image_path"] = parameters["init_img"] # TODO: Noncompliant
|
||||
rfc_dict["sampler"] = "ddim" # TODO: FIX ME WHEN IMG2IMG SUPPORTS ALL SAMPLERS
|
||||
if "init_mask" in parameters:
|
||||
rfc_dict["mask_hash"] = calculate_init_img_hash(
|
||||
parameters["init_mask"]
|
||||
) # TODO: Noncompliant
|
||||
rfc_dict["mask_image_path"] = parameters["init_mask"] # TODO: Noncompliant
|
||||
else:
|
||||
rfc_dict["type"] = "txt2img"
|
||||
|
||||
metadata["image"] = rfc_dict
|
||||
|
||||
return metadata
|
||||
|
||||
|
||||
def make_unique_init_image_filename(name):
|
||||
uuid = uuid4().hex
|
||||
split = os.path.splitext(name)
|
||||
name = f"{split[0]}.{uuid}{split[1]}"
|
||||
return name
|
||||
|
||||
|
||||
def write_log_message(message, log_path=log_path):
|
||||
"""Logs the filename and parameters used to generate or process that image to log file"""
|
||||
message = f"{message}\n"
|
||||
with open(log_path, "a", encoding="utf-8") as file:
|
||||
file.writelines(message)
|
||||
|
||||
|
||||
def save_image(
|
||||
image, command, metadata, output_dir, step_index=None, postprocessing=False
|
||||
):
|
||||
pngwriter = PngWriter(output_dir)
|
||||
prefix = pngwriter.unique_prefix()
|
||||
|
||||
seed = "unknown_seed"
|
||||
|
||||
if "image" in metadata:
|
||||
if "seed" in metadata["image"]:
|
||||
seed = metadata["image"]["seed"]
|
||||
|
||||
filename = f"{prefix}.{seed}"
|
||||
|
||||
if step_index:
|
||||
filename += f".{step_index}"
|
||||
if postprocessing:
|
||||
filename += f".postprocessed"
|
||||
|
||||
filename += ".png"
|
||||
|
||||
path = pngwriter.save_image_and_prompt_to_png(
|
||||
image=image, dream_prompt=command, metadata=metadata, name=filename
|
||||
)
|
||||
|
||||
return path
|
||||
|
||||
|
||||
def calculate_real_steps(steps, strength, has_init_image):
|
||||
return math.floor(strength * steps) if has_init_image else steps
|
||||
|
||||
|
||||
def generate_images(generation_parameters, esrgan_parameters, gfpgan_parameters):
|
||||
canceled.clear()
|
||||
|
||||
step_index = 1
|
||||
prior_variations = (
|
||||
generation_parameters["with_variations"]
|
||||
if "with_variations" in generation_parameters
|
||||
else []
|
||||
)
|
||||
"""
|
||||
If a result image is used as an init image, and then deleted, we will want to be
|
||||
able to use it as an init image in the future. Need to copy it.
|
||||
|
||||
If the init/mask image doesn't exist in the init_image_path/mask_image_path,
|
||||
make a unique filename for it and copy it there.
|
||||
"""
|
||||
if "init_img" in generation_parameters:
|
||||
filename = os.path.basename(generation_parameters["init_img"])
|
||||
if not os.path.exists(os.path.join(init_image_path, filename)):
|
||||
unique_filename = make_unique_init_image_filename(filename)
|
||||
new_path = os.path.join(init_image_path, unique_filename)
|
||||
shutil.copy(generation_parameters["init_img"], new_path)
|
||||
generation_parameters["init_img"] = new_path
|
||||
if "init_mask" in generation_parameters:
|
||||
filename = os.path.basename(generation_parameters["init_mask"])
|
||||
if not os.path.exists(os.path.join(mask_image_path, filename)):
|
||||
unique_filename = make_unique_init_image_filename(filename)
|
||||
new_path = os.path.join(init_image_path, unique_filename)
|
||||
shutil.copy(generation_parameters["init_img"], new_path)
|
||||
generation_parameters["init_mask"] = new_path
|
||||
|
||||
totalSteps = calculate_real_steps(
|
||||
steps=generation_parameters["steps"],
|
||||
strength=generation_parameters["strength"]
|
||||
if "strength" in generation_parameters
|
||||
else None,
|
||||
has_init_image="init_img" in generation_parameters,
|
||||
)
|
||||
|
||||
progress = {
|
||||
"currentStep": 1,
|
||||
"totalSteps": totalSteps,
|
||||
"currentIteration": 1,
|
||||
"totalIterations": generation_parameters["iterations"],
|
||||
"currentStatus": "Preparing",
|
||||
"isProcessing": True,
|
||||
"currentStatusHasSteps": False,
|
||||
}
|
||||
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
def image_progress(sample, step):
|
||||
if canceled.is_set():
|
||||
raise CanceledException
|
||||
|
||||
nonlocal step_index
|
||||
nonlocal generation_parameters
|
||||
nonlocal progress
|
||||
|
||||
progress["currentStep"] = step + 1
|
||||
progress["currentStatus"] = "Generating"
|
||||
progress["currentStatusHasSteps"] = True
|
||||
|
||||
if (
|
||||
generation_parameters["progress_images"]
|
||||
and step % 5 == 0
|
||||
and step < generation_parameters["steps"] - 1
|
||||
):
|
||||
image = generate.sample_to_image(sample)
|
||||
|
||||
metadata = parameters_to_generated_image_metadata(generation_parameters)
|
||||
command = parameters_to_command(generation_parameters)
|
||||
path = save_image(image, command, metadata, intermediate_path, step_index=step_index, postprocessing=False)
|
||||
|
||||
step_index += 1
|
||||
socketio.emit(
|
||||
"intermediateResult",
|
||||
{
|
||||
"url": os.path.relpath(path),
|
||||
"mtime": os.path.getmtime(path),
|
||||
"metadata": metadata,
|
||||
},
|
||||
)
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
def image_done(image, seed, first_seed):
|
||||
nonlocal generation_parameters
|
||||
nonlocal esrgan_parameters
|
||||
nonlocal gfpgan_parameters
|
||||
nonlocal progress
|
||||
|
||||
step_index = 1
|
||||
nonlocal prior_variations
|
||||
|
||||
progress["currentStatus"] = "Generation complete"
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
all_parameters = generation_parameters
|
||||
postprocessing = False
|
||||
|
||||
if (
|
||||
"variation_amount" in all_parameters
|
||||
and all_parameters["variation_amount"] > 0
|
||||
):
|
||||
first_seed = first_seed or seed
|
||||
this_variation = [[seed, all_parameters["variation_amount"]]]
|
||||
all_parameters["with_variations"] = prior_variations + this_variation
|
||||
all_parameters["seed"] = first_seed
|
||||
elif ("with_variations" in all_parameters):
|
||||
all_parameters["seed"] = first_seed
|
||||
else:
|
||||
all_parameters["seed"] = seed
|
||||
|
||||
if esrgan_parameters:
|
||||
progress["currentStatus"] = "Upscaling"
|
||||
progress["currentStatusHasSteps"] = False
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
image = esrgan.process(
|
||||
image=image,
|
||||
upsampler_scale=esrgan_parameters["level"],
|
||||
strength=esrgan_parameters["strength"],
|
||||
seed=seed,
|
||||
)
|
||||
|
||||
postprocessing = True
|
||||
all_parameters["upscale"] = [
|
||||
esrgan_parameters["level"],
|
||||
esrgan_parameters["strength"],
|
||||
]
|
||||
|
||||
if gfpgan_parameters:
|
||||
progress["currentStatus"] = "Fixing faces"
|
||||
progress["currentStatusHasSteps"] = False
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
image = gfpgan.process(
|
||||
image=image, strength=gfpgan_parameters["strength"], seed=seed
|
||||
)
|
||||
postprocessing = True
|
||||
all_parameters["gfpgan_strength"] = gfpgan_parameters["strength"]
|
||||
|
||||
progress["currentStatus"] = "Saving image"
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
metadata = parameters_to_generated_image_metadata(all_parameters)
|
||||
command = parameters_to_command(all_parameters)
|
||||
|
||||
path = save_image(
|
||||
image, command, metadata, result_path, postprocessing=postprocessing
|
||||
)
|
||||
|
||||
print(f'>> Image generated: "{path}"')
|
||||
write_log_message(f'[Generated] "{path}": {command}')
|
||||
|
||||
if progress["totalIterations"] > progress["currentIteration"]:
|
||||
progress["currentStep"] = 1
|
||||
progress["currentIteration"] += 1
|
||||
progress["currentStatus"] = "Iteration finished"
|
||||
progress["currentStatusHasSteps"] = False
|
||||
else:
|
||||
progress["currentStep"] = 0
|
||||
progress["totalSteps"] = 0
|
||||
progress["currentIteration"] = 0
|
||||
progress["totalIterations"] = 0
|
||||
progress["currentStatus"] = "Finished"
|
||||
progress["isProcessing"] = False
|
||||
|
||||
socketio.emit("progressUpdate", progress)
|
||||
eventlet.sleep(0)
|
||||
|
||||
socketio.emit(
|
||||
"generationResult",
|
||||
{
|
||||
"url": os.path.relpath(path),
|
||||
"mtime": os.path.getmtime(path),
|
||||
"metadata": metadata,
|
||||
},
|
||||
)
|
||||
eventlet.sleep(0)
|
||||
|
||||
try:
|
||||
generate.prompt2image(
|
||||
**generation_parameters,
|
||||
step_callback=image_progress,
|
||||
image_callback=image_done,
|
||||
)
|
||||
|
||||
except KeyboardInterrupt:
|
||||
raise
|
||||
except CanceledException:
|
||||
pass
|
||||
except Exception as e:
|
||||
socketio.emit("error", {"message": (str(e))})
|
||||
print("\n")
|
||||
traceback.print_exc()
|
||||
print("\n")
|
||||
|
||||
|
||||
"""
|
||||
END ADDITIONAL FUNCTIONS
|
||||
"""
|
||||
|
||||
|
||||
if __name__ == "__main__":
|
||||
print(f">> Starting server at http://{host}:{port}")
|
||||
socketio.run(app, host=host, port=port)
|
||||
18
configs/models.yaml
Normal file
@@ -0,0 +1,18 @@
|
||||
# This file describes the alternative machine learning models
|
||||
# available to the dream script.
|
||||
#
|
||||
# To add a new model, follow the examples below. Each
|
||||
# model requires a model config file, a weights file,
|
||||
# and the width and height of the images it
|
||||
# was trained on.
|
||||
|
||||
laion400m:
|
||||
config: configs/latent-diffusion/txt2img-1p4B-eval.yaml
|
||||
weights: models/ldm/text2img-large/model.ckpt
|
||||
width: 256
|
||||
height: 256
|
||||
stable-diffusion-1.4:
|
||||
config: configs/stable-diffusion/v1-inference.yaml
|
||||
weights: models/ldm/stable-diffusion-v1/model.ckpt
|
||||
width: 512
|
||||
height: 512
|
||||
@@ -73,8 +73,8 @@ model:
|
||||
data:
|
||||
target: main.DataModuleFromConfig
|
||||
params:
|
||||
batch_size: 2
|
||||
num_workers: 16
|
||||
batch_size: 1
|
||||
num_workers: 2
|
||||
wrap: false
|
||||
train:
|
||||
target: ldm.data.personalized.PersonalizedBase
|
||||
@@ -92,6 +92,9 @@ data:
|
||||
repeats: 10
|
||||
|
||||
lightning:
|
||||
modelcheckpoint:
|
||||
params:
|
||||
every_n_train_steps: 500
|
||||
callbacks:
|
||||
image_logger:
|
||||
target: main.ImageLogger
|
||||
@@ -102,4 +105,6 @@ lightning:
|
||||
|
||||
trainer:
|
||||
benchmark: True
|
||||
max_steps: 6100
|
||||
max_steps: 4000000
|
||||
# max_steps: 4000
|
||||
|
||||
|
||||
@@ -30,7 +30,7 @@ model:
|
||||
target: ldm.modules.embedding_manager.EmbeddingManager
|
||||
params:
|
||||
placeholder_strings: ["*"]
|
||||
initializer_words: ["sculpture"]
|
||||
initializer_words: ['face', 'man', 'photo', 'africanmale']
|
||||
per_image_tokens: false
|
||||
num_vectors_per_token: 1
|
||||
progressive_words: False
|
||||
|
||||
110
configs/stable-diffusion/v1-m1-finetune.yaml
Normal file
@@ -0,0 +1,110 @@
|
||||
model:
|
||||
base_learning_rate: 5.0e-03
|
||||
target: ldm.models.diffusion.ddpm.LatentDiffusion
|
||||
params:
|
||||
linear_start: 0.00085
|
||||
linear_end: 0.0120
|
||||
num_timesteps_cond: 1
|
||||
log_every_t: 200
|
||||
timesteps: 1000
|
||||
first_stage_key: image
|
||||
cond_stage_key: caption
|
||||
image_size: 64
|
||||
channels: 4
|
||||
cond_stage_trainable: true # Note: different from the one we trained before
|
||||
conditioning_key: crossattn
|
||||
monitor: val/loss_simple_ema
|
||||
scale_factor: 0.18215
|
||||
use_ema: False
|
||||
embedding_reg_weight: 0.0
|
||||
|
||||
personalization_config:
|
||||
target: ldm.modules.embedding_manager.EmbeddingManager
|
||||
params:
|
||||
placeholder_strings: ["*"]
|
||||
initializer_words: ['face', 'man', 'photo', 'africanmale']
|
||||
per_image_tokens: false
|
||||
num_vectors_per_token: 6
|
||||
progressive_words: False
|
||||
|
||||
unet_config:
|
||||
target: ldm.modules.diffusionmodules.openaimodel.UNetModel
|
||||
params:
|
||||
image_size: 32 # unused
|
||||
in_channels: 4
|
||||
out_channels: 4
|
||||
model_channels: 320
|
||||
attention_resolutions: [ 4, 2, 1 ]
|
||||
num_res_blocks: 2
|
||||
channel_mult: [ 1, 2, 4, 4 ]
|
||||
num_heads: 8
|
||||
use_spatial_transformer: True
|
||||
transformer_depth: 1
|
||||
context_dim: 768
|
||||
use_checkpoint: True
|
||||
legacy: False
|
||||
|
||||
first_stage_config:
|
||||
target: ldm.models.autoencoder.AutoencoderKL
|
||||
params:
|
||||
embed_dim: 4
|
||||
monitor: val/rec_loss
|
||||
ddconfig:
|
||||
double_z: true
|
||||
z_channels: 4
|
||||
resolution: 256
|
||||
in_channels: 3
|
||||
out_ch: 3
|
||||
ch: 128
|
||||
ch_mult:
|
||||
- 1
|
||||
- 2
|
||||
- 4
|
||||
- 4
|
||||
num_res_blocks: 2
|
||||
attn_resolutions: []
|
||||
dropout: 0.0
|
||||
lossconfig:
|
||||
target: torch.nn.Identity
|
||||
|
||||
cond_stage_config:
|
||||
target: ldm.modules.encoders.modules.FrozenCLIPEmbedder
|
||||
|
||||
data:
|
||||
target: main.DataModuleFromConfig
|
||||
params:
|
||||
batch_size: 1
|
||||
num_workers: 2
|
||||
wrap: false
|
||||
train:
|
||||
target: ldm.data.personalized.PersonalizedBase
|
||||
params:
|
||||
size: 512
|
||||
set: train
|
||||
per_image_tokens: false
|
||||
repeats: 100
|
||||
validation:
|
||||
target: ldm.data.personalized.PersonalizedBase
|
||||
params:
|
||||
size: 512
|
||||
set: val
|
||||
per_image_tokens: false
|
||||
repeats: 10
|
||||
|
||||
lightning:
|
||||
modelcheckpoint:
|
||||
params:
|
||||
every_n_train_steps: 500
|
||||
callbacks:
|
||||
image_logger:
|
||||
target: main.ImageLogger
|
||||
params:
|
||||
batch_frequency: 500
|
||||
max_images: 5
|
||||
increase_log_steps: False
|
||||
|
||||
trainer:
|
||||
benchmark: False
|
||||
max_steps: 6200
|
||||
# max_steps: 4000
|
||||
|
||||
57
docker-build/Dockerfile
Normal file
@@ -0,0 +1,57 @@
|
||||
FROM debian
|
||||
|
||||
ARG gsd
|
||||
ENV GITHUB_STABLE_DIFFUSION $gsd
|
||||
|
||||
ARG rsd
|
||||
ENV REQS $rsd
|
||||
|
||||
ARG cs
|
||||
ENV CONDA_SUBDIR $cs
|
||||
|
||||
ENV PIP_EXISTS_ACTION="w"
|
||||
|
||||
# TODO: Optimize image size
|
||||
|
||||
SHELL ["/bin/bash", "-c"]
|
||||
|
||||
WORKDIR /
|
||||
RUN apt update && apt upgrade -y \
|
||||
&& apt install -y \
|
||||
git \
|
||||
libgl1-mesa-glx \
|
||||
libglib2.0-0 \
|
||||
pip \
|
||||
python3 \
|
||||
&& git clone $GITHUB_STABLE_DIFFUSION
|
||||
|
||||
# Install Anaconda or Miniconda
|
||||
COPY anaconda.sh .
|
||||
RUN bash anaconda.sh -b -u -p /anaconda && /anaconda/bin/conda init bash
|
||||
|
||||
# SD
|
||||
WORKDIR /stable-diffusion
|
||||
RUN source ~/.bashrc \
|
||||
&& conda create -y --name ldm && conda activate ldm \
|
||||
&& conda config --env --set subdir $CONDA_SUBDIR \
|
||||
&& pip3 install -r $REQS \
|
||||
&& pip3 install basicsr facexlib realesrgan \
|
||||
&& mkdir models/ldm/stable-diffusion-v1 \
|
||||
&& ln -s "/data/sd-v1-4.ckpt" models/ldm/stable-diffusion-v1/model.ckpt
|
||||
|
||||
# Face restoreation
|
||||
# by default expected in a sibling directory to stable-diffusion
|
||||
WORKDIR /
|
||||
RUN git clone https://github.com/TencentARC/GFPGAN.git
|
||||
|
||||
WORKDIR /GFPGAN
|
||||
RUN pip3 install -r requirements.txt \
|
||||
&& python3 setup.py develop \
|
||||
&& ln -s "/data/GFPGANv1.4.pth" experiments/pretrained_models/GFPGANv1.4.pth
|
||||
|
||||
WORKDIR /stable-diffusion
|
||||
RUN python3 scripts/preload_models.py
|
||||
|
||||
WORKDIR /
|
||||
COPY entrypoint.sh .
|
||||
ENTRYPOINT ["/entrypoint.sh"]
|
||||
10
docker-build/entrypoint.sh
Executable file
@@ -0,0 +1,10 @@
|
||||
#!/bin/bash
|
||||
|
||||
cd /stable-diffusion
|
||||
|
||||
if [ $# -eq 0 ]; then
|
||||
python3 scripts/dream.py --full_precision -o /data
|
||||
# bash
|
||||
else
|
||||
python3 scripts/dream.py --full_precision -o /data "$@"
|
||||
fi
|
||||
192
docs/CHANGELOG.md
Normal file
@@ -0,0 +1,192 @@
|
||||
---
|
||||
title: Changelog
|
||||
---
|
||||
|
||||
# :octicons-log-16: **Changelog**
|
||||
|
||||
## v2.0.1 (13 October 2022)
|
||||
|
||||
- fix noisy images at high step count when using k* samplers
|
||||
- dream.py script now calls invoke.py module directly rather than
|
||||
via a new python process (which could break the environment)
|
||||
|
||||
## v2.0.0 <small>(9 October 2022)</small>
|
||||
|
||||
- `dream.py` script renamed `invoke.py`. A `dream.py` script wrapper remains
|
||||
for backward compatibility.
|
||||
- Completely new WebGUI - launch with `python3 scripts/invoke.py --web`
|
||||
- Support for <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/INPAINTING.md">inpainting</a> and <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/OUTPAINTING.md">outpainting</a>
|
||||
- img2img runs on all k* samplers
|
||||
- Support for <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/PROMPTS.md#negative-and-unconditioned-prompts">negative prompts</a>
|
||||
- Support for CodeFormer face reconstruction
|
||||
- Support for Textual Inversion on Macintoshes
|
||||
- Support in both WebGUI and CLI for <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/POSTPROCESS.md">post-processing of previously-generated images</a>
|
||||
using facial reconstruction, ESRGAN upscaling, outcropping (similar to DALL-E infinite canvas),
|
||||
and "embiggen" upscaling. See the `!fix` command.
|
||||
- New `--hires` option on `invoke>` line allows <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/CLI.m#this-is-an-example-of-txt2img">larger images to be created without duplicating elements</a>, at the cost of some performance.
|
||||
- New `--perlin` and `--threshold` options allow you to add and control variation
|
||||
during image generation (see <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/OTHER.md#thresholding-and-perlin-noise-initialization-options">Thresholding and Perlin Noise Initialization</a>
|
||||
- Extensive metadata now written into PNG files, allowing reliable regeneration of images
|
||||
and tweaking of previous settings.
|
||||
- Command-line completion in `invoke.py` now works on Windows, Linux and Mac platforms.
|
||||
- Improved <a href="https://github.com/invoke-ai/InvokeAI/blob/main/docs/features/CLI.m">command-line completion behavior</a>.
|
||||
New commands added:
|
||||
* List command-line history with `!history`
|
||||
* Search command-line history with `!search`
|
||||
* Clear history with `!clear`
|
||||
- Deprecated `--full_precision` / `-F`. Simply omit it and `invoke.py` will auto
|
||||
configure. To switch away from auto use the new flag like `--precision=float32`.
|
||||
|
||||
## v1.14 <small>(11 September 2022)</small>
|
||||
|
||||
- Memory optimizations for small-RAM cards. 512x512 now possible on 4 GB GPUs.
|
||||
- Full support for Apple hardware with M1 or M2 chips.
|
||||
- Add "seamless mode" for circular tiling of image. Generates beautiful effects.
|
||||
([prixt](https://github.com/prixt)).
|
||||
- Inpainting support.
|
||||
- Improved web server GUI.
|
||||
- Lots of code and documentation cleanups.
|
||||
|
||||
## v1.13 <small>(3 September 2022)</small>
|
||||
|
||||
- Support image variations (see [VARIATIONS](features/VARIATIONS.md)
|
||||
([Kevin Gibbons](https://github.com/bakkot) and many contributors and reviewers)
|
||||
- Supports a Google Colab notebook for a standalone server running on Google hardware
|
||||
[Arturo Mendivil](https://github.com/artmen1516)
|
||||
- WebUI supports GFPGAN/ESRGAN facial reconstruction and upscaling
|
||||
[Kevin Gibbons](https://github.com/bakkot)
|
||||
- WebUI supports incremental display of in-progress images during generation
|
||||
[Kevin Gibbons](https://github.com/bakkot)
|
||||
- A new configuration file scheme that allows new models (including upcoming
|
||||
stable-diffusion-v1.5) to be added without altering the code.
|
||||
([David Wager](https://github.com/maddavid12))
|
||||
- Can specify --grid on invoke.py command line as the default.
|
||||
- Miscellaneous internal bug and stability fixes.
|
||||
- Works on M1 Apple hardware.
|
||||
- Multiple bug fixes.
|
||||
|
||||
---
|
||||
|
||||
## v1.12 <small>(28 August 2022)</small>
|
||||
|
||||
- Improved file handling, including ability to read prompts from standard input.
|
||||
(kudos to [Yunsaki](https://github.com/yunsaki)
|
||||
- The web server is now integrated with the invoke.py script. Invoke by adding --web to
|
||||
the invoke.py command arguments.
|
||||
- Face restoration and upscaling via GFPGAN and Real-ESGAN are now automatically
|
||||
enabled if the GFPGAN directory is located as a sibling to Stable Diffusion.
|
||||
VRAM requirements are modestly reduced. Thanks to both [Blessedcoolant](https://github.com/blessedcoolant) and
|
||||
[Oceanswave](https://github.com/oceanswave) for their work on this.
|
||||
- You can now swap samplers on the invoke> command line. [Blessedcoolant](https://github.com/blessedcoolant)
|
||||
|
||||
---
|
||||
|
||||
## v1.11 <small>(26 August 2022)</small>
|
||||
|
||||
- NEW FEATURE: Support upscaling and face enhancement using the GFPGAN module. (kudos to [Oceanswave](https://github.com/Oceanswave)
|
||||
- You now can specify a seed of -1 to use the previous image's seed, -2 to use the seed for the image generated before that, etc.
|
||||
Seed memory only extends back to the previous command, but will work on all images generated with the -n# switch.
|
||||
- Variant generation support temporarily disabled pending more general solution.
|
||||
- Created a feature branch named **yunsaki-morphing-invoke** which adds experimental support for
|
||||
iteratively modifying the prompt and its parameters. Please see[ Pull Request #86](https://github.com/lstein/stable-diffusion/pull/86)
|
||||
for a synopsis of how this works. Note that when this feature is eventually added to the main branch, it will may be modified
|
||||
significantly.
|
||||
|
||||
---
|
||||
|
||||
## v1.10 <small>(25 August 2022)</small>
|
||||
|
||||
- A barebones but fully functional interactive web server for online generation of txt2img and img2img.
|
||||
|
||||
---
|
||||
|
||||
## v1.09 <small>(24 August 2022)</small>
|
||||
|
||||
- A new -v option allows you to generate multiple variants of an initial image
|
||||
in img2img mode. (kudos to [Oceanswave](https://github.com/Oceanswave). [
|
||||
See this discussion in the PR for examples and details on use](https://github.com/lstein/stable-diffusion/pull/71#issuecomment-1226700810))
|
||||
- Added ability to personalize text to image generation (kudos to [Oceanswave](https://github.com/Oceanswave) and [nicolai256](https://github.com/nicolai256))
|
||||
- Enabled all of the samplers from k_diffusion
|
||||
|
||||
---
|
||||
|
||||
## v1.08 <small>(24 August 2022)</small>
|
||||
|
||||
- Escape single quotes on the invoke> command before trying to parse. This avoids
|
||||
parse errors.
|
||||
- Removed instruction to get Python3.8 as first step in Windows install.
|
||||
Anaconda3 does it for you.
|
||||
- Added bounds checks for numeric arguments that could cause crashes.
|
||||
- Cleaned up the copyright and license agreement files.
|
||||
|
||||
---
|
||||
|
||||
## v1.07 <small>(23 August 2022)</small>
|
||||
|
||||
- Image filenames will now never fill gaps in the sequence, but will be assigned the
|
||||
next higher name in the chosen directory. This ensures that the alphabetic and chronological
|
||||
sort orders are the same.
|
||||
|
||||
---
|
||||
|
||||
## v1.06 <small>(23 August 2022)</small>
|
||||
|
||||
- Added weighted prompt support contributed by [xraxra](https://github.com/xraxra)
|
||||
- Example of using weighted prompts to tweak a demonic figure contributed by [bmaltais](https://github.com/bmaltais)
|
||||
|
||||
---
|
||||
|
||||
## v1.05 <small>(22 August 2022 - after the drop)</small>
|
||||
|
||||
- Filenames now use the following formats:
|
||||
000010.95183149.png -- Two files produced by the same command (e.g. -n2),
|
||||
000010.26742632.png -- distinguished by a different seed.
|
||||
|
||||
000011.455191342.01.png -- Two files produced by the same command using
|
||||
000011.455191342.02.png -- a batch size>1 (e.g. -b2). They have the same seed.
|
||||
|
||||
000011.4160627868.grid#1-4.png -- a grid of four images (-g); the whole grid can
|
||||
be regenerated with the indicated key
|
||||
|
||||
- It should no longer be possible for one image to overwrite another
|
||||
- You can use the "cd" and "pwd" commands at the invoke> prompt to set and retrieve
|
||||
the path of the output directory.
|
||||
|
||||
---
|
||||
|
||||
## v1.04 <small>(22 August 2022 - after the drop)</small>
|
||||
|
||||
- Updated README to reflect installation of the released weights.
|
||||
- Suppressed very noisy and inconsequential warning when loading the frozen CLIP
|
||||
tokenizer.
|
||||
|
||||
---
|
||||
|
||||
## v1.03 <small>(22 August 2022)</small>
|
||||
|
||||
- The original txt2img and img2img scripts from the CompViz repository have been moved into
|
||||
a subfolder named "orig_scripts", to reduce confusion.
|
||||
|
||||
---
|
||||
|
||||
## v1.02 <small>(21 August 2022)</small>
|
||||
|
||||
- A copy of the prompt and all of its switches and options is now stored in the corresponding
|
||||
image in a tEXt metadata field named "Dream". You can read the prompt using scripts/images2prompt.py,
|
||||
or an image editor that allows you to explore the full metadata.
|
||||
**Please run "conda env update" to load the k_lms dependencies!!**
|
||||
|
||||
---
|
||||
|
||||
## v1.01 <small>(21 August 2022)</small>
|
||||
|
||||
- added k_lms sampling.
|
||||
**Please run "conda env update" to load the k_lms dependencies!!**
|
||||
- use half precision arithmetic by default, resulting in faster execution and lower memory requirements
|
||||
Pass argument --full_precision to invoke.py to get slower but more accurate image generation
|
||||
|
||||
---
|
||||
|
||||
## Links
|
||||
|
||||
- **[Read Me](index.md)**
|
||||
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|
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|
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|
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|
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|
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|
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docs/assets/join-us-on-discord-image.png
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|
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docs/assets/logo.png
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|
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BIN
docs/assets/negative_prompt_walkthru/step1.png
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|
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docs/assets/negative_prompt_walkthru/step2.png
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|
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|
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|
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|
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|
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|
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|
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|
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|
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|
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|
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|
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|
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|
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|
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|
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BIN
docs/assets/step1.png
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|
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|
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|
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|
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|
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|
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docs/assets/truncation_comparison.jpg
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|
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|
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|
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|
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|
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docs/assets/variation_walkthru/000003.1614299449.png
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|
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|
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143
docs/features/CHANGELOG.md
Normal file
@@ -0,0 +1,143 @@
|
||||
---
|
||||
title: Changelog
|
||||
---
|
||||
|
||||
# :octicons-log-16: Changelog
|
||||
|
||||
## v1.13
|
||||
|
||||
- Supports a Google Colab notebook for a standalone server running on Google
|
||||
hardware [Arturo Mendivil](https://github.com/artmen1516)
|
||||
- WebUI supports GFPGAN/ESRGAN facial reconstruction and upscaling
|
||||
[Kevin Gibbons](https://github.com/bakkot)
|
||||
- WebUI supports incremental display of in-progress images during generation
|
||||
[Kevin Gibbons](https://github.com/bakkot)
|
||||
- Output directory can be specified on the invoke> command line.
|
||||
- The grid was displaying duplicated images when not enough images to fill the
|
||||
final row [Muhammad Usama](https://github.com/SMUsamaShah)
|
||||
- Can specify --grid on invoke.py command line as the default.
|
||||
- Miscellaneous internal bug and stability fixes.
|
||||
|
||||
---
|
||||
|
||||
## v1.12 <small>(28 August 2022)</small>
|
||||
|
||||
- Improved file handling, including ability to read prompts from standard input.
|
||||
(kudos to [Yunsaki](https://github.com/yunsaki)
|
||||
- The web server is now integrated with the invoke.py script. Invoke by adding
|
||||
--web to the invoke.py command arguments.
|
||||
- Face restoration and upscaling via GFPGAN and Real-ESGAN are now automatically
|
||||
enabled if the GFPGAN directory is located as a sibling to Stable Diffusion.
|
||||
VRAM requirements are modestly reduced. Thanks to both
|
||||
[Blessedcoolant](https://github.com/blessedcoolant) and
|
||||
[Oceanswave](https://github.com/oceanswave) for their work on this.
|
||||
- You can now swap samplers on the invoke> command line.
|
||||
[Blessedcoolant](https://github.com/blessedcoolant)
|
||||
|
||||
---
|
||||
|
||||
## v1.11 <small>(26 August 2022)</small>
|
||||
|
||||
- NEW FEATURE: Support upscaling and face enhancement using the GFPGAN module.
|
||||
(kudos to [Oceanswave](https://github.com/Oceanswave))
|
||||
- You now can specify a seed of -1 to use the previous image's seed, -2 to use
|
||||
the seed for the image generated before that, etc. Seed memory only extends
|
||||
back to the previous command, but will work on all images generated with the
|
||||
-n# switch.
|
||||
- Variant generation support temporarily disabled pending more general solution.
|
||||
- Created a feature branch named **yunsaki-morphing-invoke** which adds
|
||||
experimental support for iteratively modifying the prompt and its parameters.
|
||||
Please
|
||||
see[ Pull Request #86](https://github.com/lstein/stable-diffusion/pull/86) for
|
||||
a synopsis of how this works. Note that when this feature is eventually added
|
||||
to the main branch, it will may be modified significantly.
|
||||
|
||||
---
|
||||
|
||||
## v1.10 <small>(25 August 2022)</small>
|
||||
|
||||
- A barebones but fully functional interactive web server for online generation
|
||||
of txt2img and img2img.
|
||||
|
||||
---
|
||||
|
||||
## v1.09 <small>(24 August 2022)</small>
|
||||
|
||||
- A new -v option allows you to generate multiple variants of an initial image
|
||||
in img2img mode. (kudos to [Oceanswave](https://github.com/Oceanswave).
|
||||
- [See this discussion in the PR for examples and details on use](https://github.com/lstein/stable-diffusion/pull/71#issuecomment-1226700810))
|
||||
- Added ability to personalize text to image generation (kudos to
|
||||
[Oceanswave](https://github.com/Oceanswave) and
|
||||
[nicolai256](https://github.com/nicolai256))
|
||||
- Enabled all of the samplers from k_diffusion
|
||||
|
||||
---
|
||||
|
||||
## v1.08 <small>(24 August 2022)</small>
|
||||
|
||||
- Escape single quotes on the invoke> command before trying to parse. This avoids
|
||||
parse errors.
|
||||
- Removed instruction to get Python3.8 as first step in Windows install.
|
||||
Anaconda3 does it for you.
|
||||
- Added bounds checks for numeric arguments that could cause crashes.
|
||||
- Cleaned up the copyright and license agreement files.
|
||||
|
||||
---
|
||||
|
||||
## v1.07 <small>(23 August 2022)</small>
|
||||
|
||||
- Image filenames will now never fill gaps in the sequence, but will be assigned
|
||||
the next higher name in the chosen directory. This ensures that the alphabetic
|
||||
and chronological sort orders are the same.
|
||||
|
||||
---
|
||||
|
||||
## v1.06 <small>(23 August 2022)</small>
|
||||
|
||||
- Added weighted prompt support contributed by
|
||||
[xraxra](https://github.com/xraxra)
|
||||
- Example of using weighted prompts to tweak a demonic figure contributed by
|
||||
[bmaltais](https://github.com/bmaltais)
|
||||
|
||||
---
|
||||
|
||||
## v1.05 <small>(22 August 2022 - after the drop)</small>
|
||||
|
||||
- Filenames now use the following formats: 000010.95183149.png -- Two files
|
||||
produced by the same command (e.g. -n2), 000010.26742632.png -- distinguished
|
||||
by a different seed.
|
||||
000011.455191342.01.png -- Two files produced by the same command using
|
||||
000011.455191342.02.png -- a batch size>1 (e.g. -b2). They have the same seed.
|
||||
000011.4160627868.grid#1-4.png -- a grid of four images (-g); the whole grid
|
||||
can be regenerated with the indicated key
|
||||
|
||||
- It should no longer be possible for one image to overwrite another
|
||||
- You can use the "cd" and "pwd" commands at the invoke> prompt to set and
|
||||
retrieve the path of the output directory.
|
||||
|
||||
## v1.04 <small>(22 August 2022 - after the drop)</small>
|
||||
|
||||
- Updated README to reflect installation of the released weights.
|
||||
- Suppressed very noisy and inconsequential warning when loading the frozen CLIP
|
||||
tokenizer.
|
||||
|
||||
## v1.03 <small>(22 August 2022)</small>
|
||||
|
||||
- The original txt2img and img2img scripts from the CompViz repository have been
|
||||
moved into a subfolder named "orig_scripts", to reduce confusion.
|
||||
|
||||
## v1.02 <small>(21 August 2022)</small>
|
||||
|
||||
- A copy of the prompt and all of its switches and options is now stored in the
|
||||
corresponding image in a tEXt metadata field named "Dream". You can read the
|
||||
prompt using scripts/images2prompt.py, or an image editor that allows you to
|
||||
explore the full metadata. **Please run "conda env update -f environment.yaml"
|
||||
to load the k_lms dependencies!!**
|
||||
|
||||
## v1.01 <small>(21 August 2022)</small>
|
||||
|
||||
- added k_lms sampling. **Please run "conda env update -f environment.yaml" to
|
||||
load the k_lms dependencies!!**
|
||||
- use half precision arithmetic by default, resulting in faster execution and
|
||||
lower memory requirements Pass argument --full_precision to invoke.py to get
|
||||
slower but more accurate image generation
|
||||
359
docs/features/CLI.md
Normal file
@@ -0,0 +1,359 @@
|
||||
---
|
||||
title: CLI
|
||||
hide:
|
||||
- toc
|
||||
---
|
||||
|
||||
# :material-bash: CLI
|
||||
|
||||
## **Interactive Command Line Interface**
|
||||
|
||||
The `invoke.py` script, located in `scripts/dream.py`, provides an interactive
|
||||
interface to image generation similar to the "invoke mothership" bot that Stable
|
||||
AI provided on its Discord server.
|
||||
|
||||
Unlike the `txt2img.py` and `img2img.py` scripts provided in the original
|
||||
[CompVis/stable-diffusion](https://github.com/CompVis/stable-diffusion) source
|
||||
code repository, the time-consuming initialization of the AI model
|
||||
initialization only happens once. After that image generation from the
|
||||
command-line interface is very fast.
|
||||
|
||||
The script uses the readline library to allow for in-line editing, command
|
||||
history (++up++ and ++down++), autocompletion, and more. To help keep track of
|
||||
which prompts generated which images, the script writes a log file of image
|
||||
names and prompts to the selected output directory.
|
||||
|
||||
In addition, as of version 1.02, it also writes the prompt into the PNG file's
|
||||
metadata where it can be retrieved using `scripts/images2prompt.py`
|
||||
|
||||
The script is confirmed to work on Linux, Windows and Mac systems.
|
||||
|
||||
!!! note
|
||||
|
||||
This script runs from the command-line or can be used as a Web application. The Web GUI is
|
||||
currently rudimentary, but a much better replacement is on its way.
|
||||
|
||||
```bash
|
||||
(invokeai) ~/stable-diffusion$ python3 ./scripts/invoke.py
|
||||
* Initializing, be patient...
|
||||
Loading model from models/ldm/text2img-large/model.ckpt
|
||||
(...more initialization messages...)
|
||||
|
||||
* Initialization done! Awaiting your command...
|
||||
invoke> ashley judd riding a camel -n2 -s150
|
||||
Outputs:
|
||||
outputs/img-samples/00009.png: "ashley judd riding a camel" -n2 -s150 -S 416354203
|
||||
outputs/img-samples/00010.png: "ashley judd riding a camel" -n2 -s150 -S 1362479620
|
||||
|
||||
invoke> "there's a fly in my soup" -n6 -g
|
||||
outputs/img-samples/00011.png: "there's a fly in my soup" -n6 -g -S 2685670268
|
||||
seeds for individual rows: [2685670268, 1216708065, 2335773498, 822223658, 714542046, 3395302430]
|
||||
invoke> q
|
||||
|
||||
# this shows how to retrieve the prompt stored in the saved image's metadata
|
||||
(invokeai) ~/stable-diffusion$ python ./scripts/images2prompt.py outputs/img_samples/*.png
|
||||
00009.png: "ashley judd riding a camel" -s150 -S 416354203
|
||||
00010.png: "ashley judd riding a camel" -s150 -S 1362479620
|
||||
00011.png: "there's a fly in my soup" -n6 -g -S 2685670268
|
||||
```
|
||||
|
||||

|
||||
|
||||
The `invoke>` prompt's arguments are pretty much identical to those used in the
|
||||
Discord bot, except you don't need to type `!invoke` (it doesn't hurt if you do).
|
||||
A significant change is that creation of individual images is now the default
|
||||
unless `--grid` (`-g`) is given. A full list is given in
|
||||
[List of prompt arguments](#list-of-prompt-arguments).
|
||||
|
||||
## Arguments
|
||||
|
||||
The script itself also recognizes a series of command-line switches that will
|
||||
change important global defaults, such as the directory for image outputs and
|
||||
the location of the model weight files.
|
||||
|
||||
### List of arguments recognized at the command line
|
||||
|
||||
These command-line arguments can be passed to `invoke.py` when you first run it
|
||||
from the Windows, Mac or Linux command line. Some set defaults that can be
|
||||
overridden on a per-prompt basis (see [List of prompt arguments](#list-of-prompt-arguments). Others
|
||||
|
||||
| Argument <img width="240" align="right"/> | Shortcut <img width="100" align="right"/> | Default <img width="320" align="right"/> | Description |
|
||||
| ----------------------------------------- | ----------------------------------------- | ---------------------------------------------- | ---------------------------------------------------------------------------------------------------- |
|
||||
| `--help` | `-h` | | Print a concise help message. |
|
||||
| `--outdir <path>` | `-o<path>` | `outputs/img_samples` | Location for generated images. |
|
||||
| `--prompt_as_dir` | `-p` | `False` | Name output directories using the prompt text. |
|
||||
| `--from_file <path>` | | `None` | Read list of prompts from a file. Use `-` to read from standard input |
|
||||
| `--model <modelname>` | | `stable-diffusion-1.4` | Loads model specified in configs/models.yaml. Currently one of "stable-diffusion-1.4" or "laion400m" |
|
||||
| `--full_precision` | `-F` | `False` | Run in slower full-precision mode. Needed for Macintosh M1/M2 hardware and some older video cards. |
|
||||
| `--web` | | `False` | Start in web server mode |
|
||||
| `--host <ip addr>` | | `localhost` | Which network interface web server should listen on. Set to 0.0.0.0 to listen on any. |
|
||||
| `--port <port>` | | `9090` | Which port web server should listen for requests on. |
|
||||
| `--config <path>` | | `configs/models.yaml` | Configuration file for models and their weights. |
|
||||
| `--iterations <int>` | `-n<int>` | `1` | How many images to generate per prompt. |
|
||||
| `--grid` | `-g` | `False` | Save all image series as a grid rather than individually. |
|
||||
| `--sampler <sampler>` | `-A<sampler>` | `k_lms` | Sampler to use. Use `-h` to get list of available samplers. |
|
||||
| `--seamless` | | `False` | Create interesting effects by tiling elements of the image. |
|
||||
| `--embedding_path <path>` | | `None` | Path to pre-trained embedding manager checkpoints, for custom models |
|
||||
| `--gfpgan_dir` | | `src/gfpgan` | Path to where GFPGAN is installed. |
|
||||
| `--gfpgan_model_path` | | `experiments/pretrained_models/GFPGANv1.4.pth` | Path to GFPGAN model file, relative to `--gfpgan_dir`. |
|
||||
| `--device <device>` | `-d<device>` | `torch.cuda.current_device()` | Device to run SD on, e.g. "cuda:0" |
|
||||
| `--free_gpu_mem` | | `False` | Free GPU memory after sampling, to allow image decoding and saving in low VRAM conditions |
|
||||
| `--precision` | | `auto` | Set model precision, default is selected by device. Options: auto, float32, float16, autocast |
|
||||
|
||||
!!! warning deprecated
|
||||
|
||||
These arguments are deprecated but still work:
|
||||
|
||||
<div align="center" markdown>
|
||||
|
||||
| Argument | Shortcut | Default | Description |
|
||||
|--------------------|------------|---------------------|--------------|
|
||||
| `--weights <path>` | | `None` | Pth to weights file; use `--model stable-diffusion-1.4` instead |
|
||||
| `--laion400m` | `-l` | `False` | Use older LAION400m weights; use `--model=laion400m` instead |
|
||||
|
||||
</div>
|
||||
|
||||
!!! tip
|
||||
|
||||
On Windows systems, you may run into
|
||||
problems when passing the invoke script standard backslashed path
|
||||
names because the Python interpreter treats "\" as an escape.
|
||||
You can either double your slashes (ick): `C:\\path\\to\\my\\file`, or
|
||||
use Linux/Mac style forward slashes (better): `C:/path/to/my/file`.
|
||||
|
||||
## List of prompt arguments
|
||||
|
||||
After the invoke.py script initializes, it will present you with a
|
||||
`invoke>` prompt. Here you can enter information to generate images
|
||||
from text ([txt2img](#txt2img)), to embellish an existing image or sketch
|
||||
([img2img](#img2img)), or to selectively alter chosen regions of the image
|
||||
([inpainting](#inpainting)).
|
||||
|
||||
### txt2img
|
||||
|
||||
!!! example
|
||||
|
||||
```bash
|
||||
invoke> waterfall and rainbow -W640 -H480
|
||||
```
|
||||
|
||||
This will create the requested image with the dimensions 640 (width)
|
||||
and 480 (height).
|
||||
|
||||
Here are the invoke> command that apply to txt2img:
|
||||
|
||||
| Argument <img width="680" align="right"/> | Shortcut <img width="420" align="right"/> | Default <img width="480" align="right"/> | Description |
|
||||
|--------------------|------------|---------------------|--------------|
|
||||
| `"my prompt"` | | | Text prompt to use. The quotation marks are optional. |
|
||||
| `--width <int>` | `-W<int>` | `512` | Width of generated image |
|
||||
| `--height <int>` | `-H<int>` | `512` | Height of generated image |
|
||||
| `--iterations <int>` | `-n<int>` | `1` | How many images to generate from this prompt |
|
||||
| `--steps <int>` | `-s<int>` | `50` | How many steps of refinement to apply |
|
||||
| `--cfg_scale <float>`| `-C<float>` | `7.5` | How hard to try to match the prompt to the generated image; any number greater than 1.0 works, but the useful range is roughly 5.0 to 20.0 |
|
||||
| `--seed <int>` | `-S<int>` | `None` | Set the random seed for the next series of images. This can be used to recreate an image generated previously.|
|
||||
| `--sampler <sampler>`| `-A<sampler>`| `k_lms` | Sampler to use. Use -h to get list of available samplers. |
|
||||
| `--hires_fix` | | | Larger images often have duplication artefacts. This option suppresses duplicates by generating the image at low res, and then using img2img to increase the resolution |
|
||||
| `--grid` | `-g` | `False` | Turn on grid mode to return a single image combining all the images generated by this prompt |
|
||||
| `--individual` | `-i` | `True` | Turn off grid mode (deprecated; leave off `--grid` instead) |
|
||||
| `--outdir <path>` | `-o<path>` | `outputs/img_samples` | Temporarily change the location of these images |
|
||||
| `--seamless` | | `False` | Activate seamless tiling for interesting effects |
|
||||
| `--log_tokenization` | `-t` | `False` | Display a color-coded list of the parsed tokens derived from the prompt |
|
||||
| `--skip_normalization`| `-x` | `False` | Weighted subprompts will not be normalized. See [Weighted Prompts](./OTHER.md#weighted-prompts) |
|
||||
| `--upscale <int> <float>` | `-U <int> <float>` | `-U 1 0.75`| Upscale image by magnification factor (2, 4), and set strength of upscaling (0.0-1.0). If strength not set, will default to 0.75. |
|
||||
| `--gfpgan_strength <float>` | `-G <float>` | `-G0` | Fix faces using the GFPGAN algorithm; argument indicates how hard the algorithm should try (0.0-1.0) |
|
||||
| `--save_original` | `-save_orig`| `False` | When upscaling or fixing faces, this will cause the original image to be saved rather than replaced. |
|
||||
| `--variation <float>` |`-v<float>`| `0.0` | Add a bit of noise (0.0=none, 1.0=high) to the image in order to generate a series of variations. Usually used in combination with `-S<seed>` and `-n<int>` to generate a series a riffs on a starting image. See [Variations](./VARIATIONS.md). |
|
||||
| `--with_variations <pattern>` | `-V<pattern>`| `None` | Combine two or more variations. See [Variations](./VARIATIONS.md) for now to use this. |
|
||||
|
||||
!!! note
|
||||
|
||||
The width and height of the image must be multiples of
|
||||
64. You can provide different values, but they will be rounded down to
|
||||
the nearest multiple of 64.
|
||||
|
||||
### img2img
|
||||
|
||||
!!! example
|
||||
|
||||
```bash
|
||||
invoke> waterfall and rainbow -I./vacation-photo.png -W640 -H480 --fit
|
||||
```
|
||||
|
||||
This will modify the indicated vacation photograph by making it more
|
||||
like the prompt. Results will vary greatly depending on what is in the
|
||||
image. We also ask to `--fit` the image into a box no bigger than
|
||||
640x480. Otherwise the image size will be identical to the provided
|
||||
photo and you may run out of memory if it is large.
|
||||
|
||||
In addition to the command-line options recognized by txt2img, img2img
|
||||
accepts additional options:
|
||||
|
||||
| Argument <img width="160" align="right"/> | Shortcut | Default | Description |
|
||||
|----------------------|-------------|-----------------|--------------|
|
||||
| `--init_img <path>` | `-I<path>` | `None` | Path to the initialization image |
|
||||
| `--fit` | `-F` | `False` | Scale the image to fit into the specified -H and -W dimensions |
|
||||
| `--strength <float>` | `-s<float>` | `0.75` | How hard to try to match the prompt to the initial image. Ranges from 0.0-0.99, with higher values replacing the initial image completely.|
|
||||
|
||||
### inpainting
|
||||
|
||||
!!! example
|
||||
|
||||
```bash
|
||||
invoke> waterfall and rainbow -I./vacation-photo.png -M./vacation-mask.png -W640 -H480 --fit
|
||||
```
|
||||
|
||||
This will do the same thing as img2img, but image alterations will
|
||||
only occur within transparent areas defined by the mask file specified
|
||||
by `-M`. You may also supply just a single initial image with the areas
|
||||
to overpaint made transparent, but you must be careful not to destroy
|
||||
the pixels underneath when you create the transparent areas. See
|
||||
[Inpainting](./INPAINTING.md) for details.
|
||||
|
||||
inpainting accepts all the arguments used for txt2img and img2img, as
|
||||
well as the --mask (-M) argument:
|
||||
|
||||
| Argument <img width="100" align="right"/> | Shortcut | Default | Description |
|
||||
|--------------------|------------|---------------------|--------------|
|
||||
| `--init_mask <path>` | `-M<path>` | `None` |Path to an image the same size as the initial_image, with areas for inpainting made transparent.|
|
||||
|
||||
## Convenience commands
|
||||
|
||||
In addition to the standard image generation arguments, there are a
|
||||
series of convenience commands that begin with !:
|
||||
|
||||
### `!fix`
|
||||
|
||||
This command runs a post-processor on a previously-generated image. It
|
||||
takes a PNG filename or path and applies your choice of the `-U`, `-G`, or
|
||||
`--embiggen` switches in order to fix faces or upscale. If you provide a
|
||||
filename, the script will look for it in the current output
|
||||
directory. Otherwise you can provide a full or partial path to the
|
||||
desired file.
|
||||
|
||||
Some examples:
|
||||
|
||||
!!! example ""
|
||||
|
||||
Upscale to 4X its original size and fix faces using codeformer:
|
||||
|
||||
```bash
|
||||
invoke> !fix 0000045.4829112.png -G1 -U4 -ft codeformer
|
||||
```
|
||||
|
||||
!!! example ""
|
||||
|
||||
Use the GFPGAN algorithm to fix faces, then upscale to 3X using --embiggen:
|
||||
|
||||
```bash
|
||||
invoke> !fix 0000045.4829112.png -G0.8 -ft gfpgan
|
||||
>> fixing outputs/img-samples/0000045.4829112.png
|
||||
>> retrieved seed 4829112 and prompt "boy enjoying a banana split"
|
||||
>> GFPGAN - Restoring Faces for image seed:4829112
|
||||
Outputs:
|
||||
[1] outputs/img-samples/000017.4829112.gfpgan-00.png: !fix "outputs/img-samples/0000045.4829112.png" -s 50 -S -W 512 -H 512 -C 7.5 -A k_lms -G 0.8
|
||||
|
||||
invoke> !fix 000017.4829112.gfpgan-00.png --embiggen 3
|
||||
...lots of text...
|
||||
Outputs:
|
||||
[2] outputs/img-samples/000018.2273800735.embiggen-00.png: !fix "outputs/img-samples/000017.243781548.gfpgan-00.png" -s 50 -S 2273800735 -W 512 -H 512 -C 7.5 -A k_lms --embiggen 3.0 0.75 0.25
|
||||
```
|
||||
|
||||
### `!fetch`
|
||||
|
||||
This command retrieves the generation parameters from a previously
|
||||
generated image and either loads them into the command line. You may
|
||||
provide either the name of a file in the current output directory, or
|
||||
a full file path.
|
||||
|
||||
```bash
|
||||
invoke> !fetch 0000015.8929913.png
|
||||
# the script returns the next line, ready for editing and running:
|
||||
invoke> a fantastic alien landscape -W 576 -H 512 -s 60 -A plms -C 7.5
|
||||
```
|
||||
|
||||
Note that this command may behave unexpectedly if given a PNG file that
|
||||
was not generated by InvokeAI.
|
||||
|
||||
### `!history`
|
||||
|
||||
The invoke script keeps track of all the commands you issue during a
|
||||
session, allowing you to re-run them. On Mac and Linux systems, it
|
||||
also writes the command-line history out to disk, giving you access to
|
||||
the most recent 1000 commands issued.
|
||||
|
||||
The `!history` command will return a numbered list of all the commands
|
||||
issued during the session (Windows), or the most recent 1000 commands
|
||||
(Mac|Linux). You can then repeat a command by using the command `!NNN`,
|
||||
where "NNN" is the history line number. For example:
|
||||
|
||||
```bash
|
||||
invoke> !history
|
||||
...
|
||||
[14] happy woman sitting under tree wearing broad hat and flowing garment
|
||||
[15] beautiful woman sitting under tree wearing broad hat and flowing garment
|
||||
[18] beautiful woman sitting under tree wearing broad hat and flowing garment -v0.2 -n6
|
||||
[20] watercolor of beautiful woman sitting under tree wearing broad hat and flowing garment -v0.2 -n6 -S2878767194
|
||||
[21] surrealist painting of beautiful woman sitting under tree wearing broad hat and flowing garment -v0.2 -n6 -S2878767194
|
||||
...
|
||||
invoke> !20
|
||||
invoke> watercolor of beautiful woman sitting under tree wearing broad hat and flowing garment -v0.2 -n6 -S2878767194
|
||||
```
|
||||
|
||||
### `!search <search string>`
|
||||
|
||||
This is similar to !history but it only returns lines that contain
|
||||
`search string`. For example:
|
||||
|
||||
```bash
|
||||
invoke> !search surreal
|
||||
[21] surrealist painting of beautiful woman sitting under tree wearing broad hat and flowing garment -v0.2 -n6 -S2878767194
|
||||
```
|
||||
|
||||
### `!clear`
|
||||
|
||||
This clears the search history from memory and disk. Be advised that
|
||||
this operation is irreversible and does not issue any warnings!
|
||||
|
||||
## Command-line editing and completion
|
||||
|
||||
The command-line offers convenient history tracking, editing, and
|
||||
command completion.
|
||||
|
||||
- To scroll through previous commands and potentially edit/reuse them, use the ++up++ and ++down++ keys.
|
||||
- To edit the current command, use the ++left++ and ++right++ keys to position the cursor, and then ++backspace++, ++delete++ or insert characters.
|
||||
- To move to the very beginning of the command, type ++ctrl+a++ (or ++command+a++ on the Mac)
|
||||
- To move to the end of the command, type ++ctrl+e++.
|
||||
- To cut a section of the command, position the cursor where you want to start cutting and type ++ctrl+k++
|
||||
- To paste a cut section back in, position the cursor where you want to paste, and type ++ctrl+y++
|
||||
|
||||
Windows users can get similar, but more limited, functionality if they
|
||||
launch `invoke.py` with the `winpty` program and have the `pyreadline3`
|
||||
library installed:
|
||||
|
||||
```batch
|
||||
> winpty python scripts\invoke.py
|
||||
```
|
||||
|
||||
On the Mac and Linux platforms, when you exit invoke.py, the last 1000
|
||||
lines of your command-line history will be saved. When you restart
|
||||
`invoke.py`, you can access the saved history using the ++up++ key.
|
||||
|
||||
In addition, limited command-line completion is installed. In various
|
||||
contexts, you can start typing your command and press ++tab++. A list of
|
||||
potential completions will be presented to you. You can then type a
|
||||
little more, hit ++tab++ again, and eventually autocomplete what you want.
|
||||
|
||||
When specifying file paths using the one-letter shortcuts, the CLI
|
||||
will attempt to complete pathnames for you. This is most handy for the
|
||||
`-I` (init image) and `-M` (init mask) paths. To initiate completion, start
|
||||
the path with a slash (`/`) or `./`. For example:
|
||||
|
||||
```bash
|
||||
invoke> zebra with a mustache -I./test-pictures<TAB>
|
||||
-I./test-pictures/Lincoln-and-Parrot.png -I./test-pictures/zebra.jpg -I./test-pictures/madonna.png
|
||||
-I./test-pictures/bad-sketch.png -I./test-pictures/man_with_eagle/
|
||||
```
|
||||
|
||||
You can then type ++z++, hit ++tab++ again, and it will autofill to `zebra.jpg`.
|
||||
|
||||
More text completion features (such as autocompleting seeds) are on their way.
|
||||
158
docs/features/EMBIGGEN.md
Normal file
@@ -0,0 +1,158 @@
|
||||
---
|
||||
title: Embiggen
|
||||
---
|
||||
|
||||
# :material-loupe: Embiggen
|
||||
|
||||
**upscale your images on limited memory machines**
|
||||
|
||||
GFPGAN and Real-ESRGAN are both memory intensive. In order to avoid
|
||||
crashes and memory overloads during the Stable Diffusion process,
|
||||
these effects are applied after Stable Diffusion has completed its
|
||||
work.
|
||||
|
||||
In single image generations, you will see the output right away but
|
||||
when you are using multiple iterations, the images will first be
|
||||
generated and then upscaled and face restored after that process is
|
||||
complete. While the image generation is taking place, you will still
|
||||
be able to preview the base images.
|
||||
|
||||
If you wish to stop during the image generation but want to upscale or
|
||||
face restore a particular generated image, pass it again with the same
|
||||
prompt and generated seed along with the `-U` and `-G` prompt
|
||||
arguments to perform those actions.
|
||||
|
||||
## Embiggen
|
||||
|
||||
If you wanted to be able to do more (pixels) without running out of VRAM,
|
||||
or you want to upscale with details that couldn't possibly appear
|
||||
without the context of a prompt, this is the feature to try out.
|
||||
|
||||
Embiggen automates the process of taking an init image, upscaling it,
|
||||
cutting it into smaller tiles that slightly overlap, running all the
|
||||
tiles through img2img to refine details with respect to the prompt,
|
||||
and "stitching" the tiles back together into a cohesive image.
|
||||
|
||||
It automatically computes how many tiles are needed, and so it can be fed
|
||||
*ANY* size init image and perform Img2Img on it (though it will be run only
|
||||
one tile at a time, which can cause problems, see the Note at the end).
|
||||
|
||||
If you're familiar with "GoBig" (ala [progrock-stable](https://github.com/lowfuel/progrock-stable))
|
||||
it's similar to that, except it can work up to an arbitrarily large size
|
||||
(instead of just 2x), with tile overlaps configurable as a ratio, and
|
||||
has extra logic to re-run any number of the tile sub-sections of the image
|
||||
if for example a small part of a huge run got messed up.
|
||||
|
||||
### Usage
|
||||
|
||||
`-embiggen <scaling_factor> <esrgan_strength> <overlap_ratio OR overlap_pixels>`
|
||||
|
||||
Takes a scaling factor relative to the size of the `--init_img` (`-I`), followed by
|
||||
ESRGAN upscaling strength (0 - 1.0), followed by minimum amount of overlap
|
||||
between tiles as a decimal ratio (0 - 1.0) *OR* a number of pixels.
|
||||
|
||||
The scaling factor is how much larger than the `--init_img` the output
|
||||
should be, and will multiply both x and y axis, so an image that is a
|
||||
scaling factor of 3.0 has 3*3= 9 times as many pixels, and will take
|
||||
(at least) 9 times as long (see overlap for why it might be
|
||||
longer). If the `--init_img` is already the right size `-embiggen 1`,
|
||||
and it can also be less than one if the init_img is too big.
|
||||
|
||||
Esrgan_strength defaults to 0.75, and the overlap_ratio defaults to
|
||||
0.25, both are optional.
|
||||
|
||||
Unlike Img2Img, the `--width` (`-W`) and `--height` (`-H`) arguments
|
||||
do not control the size of the image as a whole, but the size of the
|
||||
tiles used to Embiggen the image.
|
||||
|
||||
ESRGAN is used to upscale the `--init_img` prior to cutting it into
|
||||
tiles/pieces to run through img2img and then stitch back
|
||||
together. Embiggen can be run without ESRGAN; just set the strength to
|
||||
zero (e.g. `-embiggen 1.75 0`). The output of Embiggen can also be
|
||||
upscaled after it's finished (`-U`).
|
||||
|
||||
The overlap is the minimum that tiles will overlap with adjacent
|
||||
tiles, specified as either a ratio or a number of pixels. How much the
|
||||
tiles overlap determines the likelihood the tiling will be noticable,
|
||||
really small overlaps (e.g. a couple of pixels) may produce noticeable
|
||||
grid-like fuzzy distortions in the final stitched image. Though, as
|
||||
the overlapping space doesn't contribute to making the image bigger,
|
||||
and the larger the overlap the more tiles (and the more time) it will
|
||||
take to finish.
|
||||
|
||||
Because the overlapping parts of tiles don't "contribute" to
|
||||
increasing size, every tile after the first in a row or column
|
||||
effectively only covers an extra `1 - overlap_ratio` on each axis. If
|
||||
the input/`--init_img` is same size as a tile, the ideal (for time)
|
||||
scaling factors with the default overlap (0.25) are 1.75, 2.5, 3.25,
|
||||
4.0 etc..
|
||||
|
||||
`-embiggen_tiles <spaced list of tiles>`
|
||||
|
||||
An advanced usage useful if you only want to alter parts of the image
|
||||
while running Embiggen. It takes a list of tiles by number to run and
|
||||
replace onto the initial image e.g. `1 3 5`. It's useful for either
|
||||
fixing problem spots from a previous Embiggen run, or selectively
|
||||
altering the prompt for sections of an image - for creative or
|
||||
coherency reasons.
|
||||
|
||||
Tiles are numbered starting with one, and left-to-right,
|
||||
top-to-bottom. So, if you are generating a 3x3 tiled image, the
|
||||
middle row would be `4 5 6`.
|
||||
|
||||
### Examples
|
||||
|
||||
!!! example ""
|
||||
|
||||
Running Embiggen with 512x512 tiles on an existing image, scaling up by a factor of 2.5x;
|
||||
and doing the same again (default ESRGAN strength is 0.75, default overlap between tiles is 0.25):
|
||||
|
||||
```bash
|
||||
invoke > a photo of a forest at sunset -s 100 -W 512 -H 512 -I outputs/forest.png -f 0.4 -embiggen 2.5
|
||||
invoke > a photo of a forest at sunset -s 100 -W 512 -H 512 -I outputs/forest.png -f 0.4 -embiggen 2.5 0.75 0.25
|
||||
```
|
||||
|
||||
If your starting image was also 512x512 this should have taken 9 tiles.
|
||||
|
||||
!!! example ""
|
||||
|
||||
If there weren't enough clouds in the sky of that forest you just made
|
||||
(and that image is about 1280 pixels (512*2.5) wide A.K.A. three
|
||||
512x512 tiles with 0.25 overlaps wide) we can replace that top row of
|
||||
tiles:
|
||||
|
||||
```bash
|
||||
invoke> a photo of puffy clouds over a forest at sunset -s 100 -W 512 -H 512 -I outputs/000002.seed.png -f 0.5 -embiggen_tiles 1 2 3
|
||||
```
|
||||
|
||||
## Fixing Previously-Generated Images
|
||||
|
||||
It is easy to apply embiggen to any previously-generated file without having to
|
||||
look up the original prompt and provide an initial image. Just use the
|
||||
syntax `!fix path/to/file.png <embiggen>`. For example, you can rewrite the
|
||||
previous command to look like this:
|
||||
|
||||
```bash
|
||||
invoke> !fix ./outputs/000002.seed.png -embiggen_tiles 1 2 3
|
||||
```
|
||||
|
||||
A new file named `000002.seed.fixed.png` will be created in the output directory. Note that
|
||||
the `!fix` command does not replace the original file, unlike the behavior at generate time.
|
||||
You do not need to provide the prompt, and `!fix` automatically selects a good strength for
|
||||
embiggen-ing.
|
||||
|
||||
!!! note
|
||||
|
||||
Because the same prompt is used on all the tiled images, and the model
|
||||
doesn't have the context of anything outside the tile being run - it
|
||||
can end up creating repeated pattern (also called 'motifs') across all
|
||||
the tiles based on that prompt. The best way to combat this is
|
||||
lowering the `--strength` (`-f`) to stay more true to the init image,
|
||||
and increasing the number of steps so there is more compute-time to
|
||||
create the detail. Anecdotally `--strength` 0.35-0.45 works pretty
|
||||
well on most things. It may also work great in some examples even with
|
||||
the `--strength` set high for patterns, landscapes, or subjects that
|
||||
are more abstract. Because this is (relatively) fast, you can also
|
||||
preserve the best parts from each.
|
||||
|
||||
Author: [Travco](https://github.com/travco)
|
||||
178
docs/features/IMG2IMG.md
Normal file
@@ -0,0 +1,178 @@
|
||||
---
|
||||
title: Image-to-Image
|
||||
---
|
||||
|
||||
# :material-image-multiple: Image-to-Image
|
||||
|
||||
## `img2img`
|
||||
|
||||
This script also provides an `img2img` feature that lets you seed your creations with an initial
|
||||
drawing or photo. This is a really cool feature that tells stable diffusion to build the prompt on
|
||||
top of the image you provide, preserving the original's basic shape and layout. To use it, provide
|
||||
the `--init_img` option as shown here:
|
||||
|
||||
```commandline
|
||||
tree on a hill with a river, nature photograph, national geographic -I./test-pictures/tree-and-river-sketch.png -f 0.85
|
||||
```
|
||||
|
||||
This will take the original image shown here:
|
||||
|
||||
<div align="center" markdown>
|
||||
<img src="https://user-images.githubusercontent.com/50542132/193946000-c42a96d8-5a74-4f8a-b4c3-5213e6cadcce.png" width=350>
|
||||
</div>
|
||||
|
||||
and generate a new image based on it as shown here:
|
||||
|
||||
<div align="center" markdown>
|
||||
<img src="https://user-images.githubusercontent.com/111189/194135515-53d4c060-e994-4016-8121-7c685e281ac9.png" width=350>
|
||||
</div>
|
||||
|
||||
The `--init_img` (`-I`) option gives the path to the seed picture. `--strength` (`-f`) controls how much
|
||||
the original will be modified, ranging from `0.0` (keep the original intact), to `1.0` (ignore the
|
||||
original completely). The default is `0.75`, and ranges from `0.25-0.90` give interesting results.
|
||||
Other relevant options include `-C` (classification free guidance scale), and `-s` (steps). Unlike `txt2img`,
|
||||
adding steps will continuously change the resulting image and it will not converge.
|
||||
|
||||
You may also pass a `-v<variation_amount>` option to generate `-n<iterations>` count variants on
|
||||
the original image. This is done by passing the first generated image
|
||||
back into img2img the requested number of times. It generates
|
||||
interesting variants.
|
||||
|
||||
Note that the prompt makes a big difference. For example, this slight variation on the prompt produces
|
||||
a very different image:
|
||||
|
||||
`photograph of a tree on a hill with a river`
|
||||
|
||||
<div align="center" markdown>
|
||||
<img src="https://user-images.githubusercontent.com/111189/194135220-16b62181-b60c-4248-8989-4834a8fd7fbd.png" width=350>
|
||||
</div>
|
||||
|
||||
!!! tip
|
||||
|
||||
When designing prompts, think about how the images scraped from the internet were captioned. Very few photographs will
|
||||
be labeled "photograph" or "photorealistic." They will, however, be captioned with the publication, photographer, camera
|
||||
model, or film settings.
|
||||
|
||||
If the initial image contains transparent regions, then Stable Diffusion will only draw within the
|
||||
transparent regions, a process called [`inpainting`](./INPAINTING.md#creating-transparent-regions-for-inpainting). However, for this to work correctly, the color
|
||||
information underneath the transparent needs to be preserved, not erased.
|
||||
|
||||
!!! warning
|
||||
|
||||
`img2img` does not work properly on initial images smaller than 512x512. Please scale your
|
||||
image to at least 512x512 before using it. Larger images are not a problem, but may run out of VRAM on your
|
||||
GPU card.
|
||||
|
||||
To fix this, use the `--fit` option, which downscales the initial image to fit within the box specified
|
||||
by width x height:
|
||||
|
||||
```bash
|
||||
invoke> "tree on a hill with a river, national geographic" -I./test-pictures/big-sketch.png -H512 -W512 --fit
|
||||
```
|
||||
|
||||
## How does it actually work, though?
|
||||
|
||||
The main difference between `img2img` and `prompt2img` is the starting point. While `prompt2img` always starts with pure
|
||||
gaussian noise and progressively refines it over the requested number of steps, `img2img` skips some of these earlier steps
|
||||
(how many it skips is indirectly controlled by the `--strength` parameter), and uses instead your initial image mixed with gaussian noise as the starting image.
|
||||
|
||||
**Let's start** by thinking about vanilla `prompt2img`, just generating an image from a prompt. If the step count is 10, then the "latent space" (Stable Diffusion's internal representation of the image) for the prompt "fire" with seed `1592514025` develops something like this:
|
||||
|
||||
```bash
|
||||
invoke> "fire" -s10 -W384 -H384 -S1592514025
|
||||
```
|
||||
|
||||
<div align="center" markdown>
|
||||

|
||||
</div>
|
||||
|
||||
Put simply: starting from a frame of fuzz/static, SD finds details in each frame that it thinks look like "fire" and brings them a little bit more into focus, gradually scrubbing out the fuzz until a clear image remains.
|
||||
|
||||
**When you use `img2img`** some of the earlier steps are cut, and instead an initial image of your choice is used. But because of how the maths behind Stable Diffusion works, this image needs to be mixed with just the right amount of noise (fuzz/static) for where it is being inserted. This is where the strength parameter comes in. Depending on the set strength, your image will be inserted into the sequence at the appropriate point, with just the right amount of noise.
|
||||
|
||||
### A concrete example
|
||||
|
||||
I want SD to draw a fire based on this hand-drawn image:
|
||||
|
||||
<div align="center" markdown>
|
||||

|
||||
</div>
|
||||
|
||||
Let's only do 10 steps, to make it easier to see what's happening. If strength is `0.7`, this is what the internal steps the algorithm has to take will look like:
|
||||
|
||||
<div align="center" markdown>
|
||||

|
||||
</div>
|
||||
|
||||
With strength `0.4`, the steps look more like this:
|
||||
|
||||
<div align="center" markdown>
|
||||

|
||||
</div>
|
||||
|
||||
Notice how much more fuzzy the starting image is for strength `0.7` compared to `0.4`, and notice also how much longer the sequence is with `0.7`:
|
||||
|
||||
| | strength = 0.7 | strength = 0.4 |
|
||||
| -- | :--: | :--: |
|
||||
| initial image that SD sees |  |  |
|
||||
| steps argument to `dream>` | `-S10` | `-S10` |
|
||||
| steps actually taken | 7 | 4 |
|
||||
| latent space at each step |  |  |
|
||||
| output |  |  |
|
||||
|
||||
Both of the outputs look kind of like what I was thinking of. With the strength higher, my input becomes more vague, *and* Stable Diffusion has more steps to refine its output. But it's not really making what I want, which is a picture of cheery open fire. With the strength lower, my input is more clear, *but* Stable Diffusion has less chance to refine itself, so the result ends up inheriting all the problems of my bad drawing.
|
||||
|
||||
If you want to try this out yourself, all of these are using a seed of `1592514025` with a width/height of `384`, step count `10`, the default sampler (`k_lms`), and the single-word prompt `"fire"`:
|
||||
|
||||
```bash
|
||||
invoke> "fire" -s10 -W384 -H384 -S1592514025 -I /tmp/fire-drawing.png --strength 0.7
|
||||
```
|
||||
|
||||
The code for rendering intermediates is on my (damian0815's) branch [document-img2img](https://github.com/damian0815/InvokeAI/tree/document-img2img) - run `invoke.py` and check your `outputs/img-samples/intermediates` folder while generating an image.
|
||||
|
||||
### Compensating for the reduced step count
|
||||
|
||||
After putting this guide together I was curious to see how the difference would be if I increased the step count to compensate, so that SD could have the same amount of steps to develop the image regardless of the strength. So I ran the generation again using the same seed, but this time adapting the step count to give each generation 20 steps.
|
||||
|
||||
Here's strength `0.4` (note step count `50`, which is `20 ÷ 0.4` to make sure SD does `20` steps from my image):
|
||||
|
||||
```bash
|
||||
invoke> "fire" -s50 -W384 -H384 -S1592514025 -I /tmp/fire-drawing.png -f 0.4
|
||||
```
|
||||
|
||||
<div align="center" markdown>
|
||||

|
||||
</div>
|
||||
|
||||
and here is strength `0.7` (note step count `30`, which is roughly `20 ÷ 0.7` to make sure SD does `20` steps from my image):
|
||||
|
||||
```bash
|
||||
invoke> "fire" -s30 -W384 -H384 -S1592514025 -I /tmp/fire-drawing.png -f 0.7
|
||||
```
|
||||
|
||||
<div align="center" markdown>
|
||||

|
||||
</div>
|
||||
|
||||
In both cases the image is nice and clean and "finished", but because at strength `0.7` Stable Diffusion has been give so much more freedom to improve on my badly-drawn flames, they've come out looking much better. You can really see the difference when looking at the latent steps. There's more noise on the first image with strength `0.7`:
|
||||
|
||||

|
||||
|
||||
than there is for strength `0.4`:
|
||||
|
||||

|
||||
|
||||
and that extra noise gives the algorithm more choices when it is evaluating how to denoise any particular pixel in the image.
|
||||
|
||||
Unfortunately, it seems that `img2img` is very sensitive to the step count. Here's strength `0.7` with a step count of `29` (SD did 19 steps from my image):
|
||||
|
||||
<div align="center" markdown>
|
||||

|
||||
</div>
|
||||
|
||||
By comparing the latents we can sort of see that something got interpreted differently enough on the third or fourth step to lead to a rather different interpretation of the flames.
|
||||
|
||||

|
||||

|
||||
|
||||
This is the result of a difference in the de-noising "schedule" - basically the noise has to be cleaned by a certain degree each step or the model won't "converge" on the image properly (see [stable diffusion blog](https://huggingface.co/blog/stable_diffusion) for more about that). A different step count means a different schedule, which means things get interpreted slightly differently at every step.
|
||||