mirror of
https://github.com/nod-ai/SHARK-Studio.git
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Add Stable Diffusion Img2Img model script
This commit is contained in:
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# Stable Diffusion Img2Img model
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## Installation
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<details>
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<summary>Installation (Linux)</summary>
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### Activate shark.venv Virtual Environment
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```shell
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source shark.venv/bin/activate
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# Some older pip installs may not be able to handle the recent PyTorch deps
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python -m pip install --upgrade pip
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```
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### Install dependencies
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# Run the setup.sh script
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```shell
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./setup.sh
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```
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### Run the Stable diffusion Img2Img model
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To run the model with the default set of images and params, run:
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```shell
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python stable_diffusion_img2img.py
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```
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To run the model with your set of images, and parameters you need to specify the following params:
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1.) Input images directory with the arg `--input_dir` containing 3-5 images.
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2.) What to teach the model? Using the arg `--what_to_teach`, allowed values are `object` or `style`.
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3.) Placeholder token using the arg `--placeholder_token`, that represents your new concept. It should be passed with the opening and closing angle brackets. For ex: token is `cat-toy`, it should be passed as `<cat-toy>`.
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4.) Initializer token using the arg `--initializer_token`, which summarise what is your new concept.
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For the result, you need to pass the text prompt with the arg: `--prompt`. The prompt string should contain a "*s" in it, which will be replaced by the placeholder token during the inference.
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By default the result images will go into the `sd_result` dir. To specify your output dir use the arg: `--output_dir`.
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The default value of max_training_steps is `3000`, which takes some hours to complete. You can pass the smaller value with the arg `--training_steps`. Specify the number of images to be sampled for the result with the `--num_inference_samples` arg.
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@@ -0,0 +1,25 @@
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#!/bin/bash
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TD="$(cd $(dirname $0) && pwd)"
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if [ -z "$PYTHON" ]; then
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PYTHON="$(which python3)"
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fi
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function die() {
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echo "Error executing command: $*"
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exit 1
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}
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PYTHON_VERSION_X_Y=`${PYTHON} -c 'import sys; version=sys.version_info[:2]; print("{0}.{1}".format(*version))'`
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echo "Python: $PYTHON"
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echo "Python version: $PYTHON_VERSION_X_Y"
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mkdir input_images
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wget https://huggingface.co/datasets/valhalla/images/resolve/main/2.jpeg -P input_images/
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wget https://huggingface.co/datasets/valhalla/images/resolve/main/3.jpeg -P input_images/
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wget https://huggingface.co/datasets/valhalla/images/resolve/main/5.jpeg -P input_images/
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wget https://huggingface.co/datasets/valhalla/images/resolve/main/6.jpeg -P input_images/
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pip install diffusers["training"]==0.4.1 transformers ftfy opencv-python
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@@ -0,0 +1,597 @@
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# Textual-inversion fine-tuning for Stable Diffusion using diffusers
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# This script shows how to "teach" Stable Diffusion a new concept via
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# textual-inversion using 🤗 Hugging Face [🧨 Diffusers library](https://github.com/huggingface/diffusers).
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# By using just 3-5 images you can teach new concepts to Stable Diffusion
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# and personalize the model on your own images.
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import argparse
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import itertools
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import math
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import os
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import random
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import cv2
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import numpy as np
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import torch
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import torch.nn.functional as F
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import torch.utils.checkpoint
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from torch.utils.data import Dataset
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import PIL
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from accelerate import Accelerator
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from accelerate.logging import get_logger
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from accelerate.utils import set_seed
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from diffusers import (
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AutoencoderKL,
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DDPMScheduler,
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PNDMScheduler,
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StableDiffusionPipeline,
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UNet2DConditionModel,
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)
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from diffusers.hub_utils import init_git_repo, push_to_hub
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from diffusers.optimization import get_scheduler
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from diffusers.pipelines.stable_diffusion import StableDiffusionSafetyChecker
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from PIL import Image
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from torchvision import transforms
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from tqdm.auto import tqdm
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from transformers import CLIPFeatureExtractor, CLIPTextModel, CLIPTokenizer
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YOUR_TOKEN = "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"
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p = argparse.ArgumentParser(
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description=__doc__, formatter_class=argparse.ArgumentDefaultsHelpFormatter
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)
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p.add_argument(
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"--input_dir",
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type=str,
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default="input_images/",
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help="the directory contains the images used for fine tuning",
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)
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p.add_argument(
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"--output_dir",
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type=str,
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default="sd_result",
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help="the directory contains the images used for fine tuning",
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)
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p.add_argument(
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"--training_steps",
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type=int,
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default=3000,
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help="the maximum number of training steps",
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)
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p.add_argument("--seed", type=int, default=42, help="the random seed")
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p.add_argument(
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"--what_to_teach",
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type=str,
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choices=["object", "style"],
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default="object",
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help="what is it that you are teaching?",
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)
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p.add_argument(
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"--placeholder_token",
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type=str,
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default="<cat-toy>",
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help="It is the token you are going to use to represent your new concept",
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)
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p.add_argument(
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"--initializer_token",
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type=str,
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default="toy",
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help="It is a word that can summarise what is your new concept",
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)
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p.add_argument(
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"--inference_steps",
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type=int,
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default=50,
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help="the number of steps for inference",
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)
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p.add_argument(
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"--num_inference_samples",
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type=int,
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default=4,
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help="the number of samples for inference",
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)
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p.add_argument(
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"--prompt",
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type=str,
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default="a grafitti in a wall with a *s on it",
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help="the text prompt to use",
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)
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args = p.parse_args()
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if "*s" not in args.prompt:
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raise ValueError(
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f'The prompt should have a "*s" which will be replaced by a placeholder token.'
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)
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prompt1, prompt2 = args.prompt.split("*s")
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args.prompt = prompt1 + args.placeholder_token + prompt2
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pretrained_model_name_or_path = "CompVis/stable-diffusion-v1-4"
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# Load input images.
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images = []
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for filename in os.listdir(args.input_dir):
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img = cv2.imread(os.path.join(args.input_dir, filename))
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if img is not None:
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images.append(img)
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# Setup the prompt templates for training
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imagenet_templates_small = [
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"a photo of a {}",
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"a rendering of a {}",
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"a cropped photo of the {}",
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"the photo of a {}",
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"a photo of a clean {}",
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"a photo of a dirty {}",
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"a dark photo of the {}",
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"a photo of my {}",
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"a photo of the cool {}",
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"a close-up photo of a {}",
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"a bright photo of the {}",
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"a cropped photo of a {}",
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"a photo of the {}",
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"a good photo of the {}",
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"a photo of one {}",
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"a close-up photo of the {}",
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"a rendition of the {}",
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"a photo of the clean {}",
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"a rendition of a {}",
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"a photo of a nice {}",
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"a good photo of a {}",
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"a photo of the nice {}",
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"a photo of the small {}",
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"a photo of the weird {}",
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"a photo of the large {}",
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"a photo of a cool {}",
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"a photo of a small {}",
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]
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imagenet_style_templates_small = [
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"a painting in the style of {}",
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"a rendering in the style of {}",
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"a cropped painting in the style of {}",
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"the painting in the style of {}",
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"a clean painting in the style of {}",
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"a dirty painting in the style of {}",
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"a dark painting in the style of {}",
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"a picture in the style of {}",
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"a cool painting in the style of {}",
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"a close-up painting in the style of {}",
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"a bright painting in the style of {}",
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"a cropped painting in the style of {}",
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"a good painting in the style of {}",
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"a close-up painting in the style of {}",
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"a rendition in the style of {}",
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"a nice painting in the style of {}",
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"a small painting in the style of {}",
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"a weird painting in the style of {}",
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"a large painting in the style of {}",
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]
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# Setup the dataset
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class TextualInversionDataset(Dataset):
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def __init__(
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self,
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data_root,
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tokenizer,
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learnable_property="object", # [object, style]
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size=512,
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repeats=100,
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interpolation="bicubic",
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flip_p=0.5,
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set="train",
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placeholder_token="*",
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center_crop=False,
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):
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self.data_root = data_root
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self.tokenizer = tokenizer
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self.learnable_property = learnable_property
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self.size = size
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self.placeholder_token = placeholder_token
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self.center_crop = center_crop
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self.flip_p = flip_p
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self.image_paths = [
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os.path.join(self.data_root, file_path)
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for file_path in os.listdir(self.data_root)
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]
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self.num_images = len(self.image_paths)
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self._length = self.num_images
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if set == "train":
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self._length = self.num_images * repeats
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self.interpolation = {
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"linear": PIL.Image.LINEAR,
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"bilinear": PIL.Image.BILINEAR,
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"bicubic": PIL.Image.BICUBIC,
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"lanczos": PIL.Image.LANCZOS,
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}[interpolation]
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self.templates = (
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imagenet_style_templates_small
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if learnable_property == "style"
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else imagenet_templates_small
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)
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self.flip_transform = transforms.RandomHorizontalFlip(p=self.flip_p)
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def __len__(self):
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return self._length
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def __getitem__(self, i):
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example = {}
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image = Image.open(self.image_paths[i % self.num_images])
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if not image.mode == "RGB":
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image = image.convert("RGB")
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placeholder_string = self.placeholder_token
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text = random.choice(self.templates).format(placeholder_string)
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example["input_ids"] = self.tokenizer(
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text,
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padding="max_length",
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truncation=True,
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max_length=self.tokenizer.model_max_length,
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return_tensors="pt",
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).input_ids[0]
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# default to score-sde preprocessing
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img = np.array(image).astype(np.uint8)
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if self.center_crop:
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crop = min(img.shape[0], img.shape[1])
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h, w, = (
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img.shape[0],
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img.shape[1],
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)
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img = img[
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(h - crop) // 2 : (h + crop) // 2,
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(w - crop) // 2 : (w + crop) // 2,
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]
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image = Image.fromarray(img)
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image = image.resize(
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(self.size, self.size), resample=self.interpolation
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)
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image = self.flip_transform(image)
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image = np.array(image).astype(np.uint8)
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image = (image / 127.5 - 1.0).astype(np.float32)
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example["pixel_values"] = torch.from_numpy(image).permute(2, 0, 1)
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return example
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# Setting up the model
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# Load the tokenizer and add the placeholder token as a additional special token.
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# Please read and if you agree accept the LICENSE
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# [here](https://huggingface.co/CompVis/stable-diffusion-v1-4) if you see an error
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tokenizer = CLIPTokenizer.from_pretrained(
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pretrained_model_name_or_path,
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subfolder="tokenizer",
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use_auth_token=YOUR_TOKEN,
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)
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# Add the placeholder token in tokenizer
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num_added_tokens = tokenizer.add_tokens(args.placeholder_token)
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if num_added_tokens == 0:
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raise ValueError(
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f"The tokenizer already contains the token {args.placeholder_token}. Please pass a different"
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" `placeholder_token` that is not already in the tokenizer."
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)
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# Get token ids for our placeholder and initializer token.
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# This code block will complain if initializer string is not a single token
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# Convert the initializer_token, placeholder_token to ids
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token_ids = tokenizer.encode(args.initializer_token, add_special_tokens=False)
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# Check if initializer_token is a single token or a sequence of tokens
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if len(token_ids) > 1:
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raise ValueError("The initializer token must be a single token.")
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initializer_token_id = token_ids[0]
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placeholder_token_id = tokenizer.convert_tokens_to_ids(args.placeholder_token)
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# Load the Stable Diffusion model
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# Load models and create wrapper for stable diffusion
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text_encoder = CLIPTextModel.from_pretrained(
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pretrained_model_name_or_path,
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subfolder="text_encoder",
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use_auth_token=YOUR_TOKEN,
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)
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vae = AutoencoderKL.from_pretrained(
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pretrained_model_name_or_path,
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subfolder="vae",
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use_auth_token=YOUR_TOKEN,
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)
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unet = UNet2DConditionModel.from_pretrained(
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pretrained_model_name_or_path,
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subfolder="unet",
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use_auth_token=YOUR_TOKEN,
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)
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# We have added the `placeholder_token` in the `tokenizer` so we resize the token embeddings here,
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# this will a new embedding vector in the token embeddings for our `placeholder_token`
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text_encoder.resize_token_embeddings(len(tokenizer))
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# Initialise the newly added placeholder token with the embeddings of the initializer token
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token_embeds = text_encoder.get_input_embeddings().weight.data
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token_embeds[placeholder_token_id] = token_embeds[initializer_token_id]
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# In Textual-Inversion we only train the newly added embedding vector,
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# so lets freeze rest of the model parameters here.
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def freeze_params(params):
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for param in params:
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param.requires_grad = False
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# Freeze vae and unet
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freeze_params(vae.parameters())
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freeze_params(unet.parameters())
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# Freeze all parameters except for the token embeddings in text encoder
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params_to_freeze = itertools.chain(
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text_encoder.text_model.encoder.parameters(),
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text_encoder.text_model.final_layer_norm.parameters(),
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text_encoder.text_model.embeddings.position_embedding.parameters(),
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)
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freeze_params(params_to_freeze)
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# Creating our training data
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train_dataset = TextualInversionDataset(
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data_root=args.input_dir,
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tokenizer=tokenizer,
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size=512,
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placeholder_token=args.placeholder_token,
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repeats=100,
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learnable_property=args.what_to_teach, # Option selected above between object and style
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center_crop=False,
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set="train",
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)
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def create_dataloader(train_batch_size=1):
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return torch.utils.data.DataLoader(
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train_dataset, batch_size=train_batch_size, shuffle=True
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)
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# Create noise_scheduler for training.
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noise_scheduler = DDPMScheduler(
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beta_start=0.00085,
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beta_end=0.012,
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beta_schedule="scaled_linear",
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num_train_timesteps=1000,
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tensor_format="pt",
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)
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# Define hyperparameters for our training
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hyperparameters = {
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"learning_rate": 5e-04,
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"scale_lr": True,
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"max_train_steps": args.training_steps,
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"train_batch_size": 1,
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"gradient_accumulation_steps": 4,
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"seed": args.seed,
|
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"output_dir": "sd-concept-output",
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}
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def training_function(text_encoder, vae, unet):
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logger = get_logger(__name__)
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train_batch_size = hyperparameters["train_batch_size"]
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gradient_accumulation_steps = hyperparameters[
|
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"gradient_accumulation_steps"
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]
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learning_rate = hyperparameters["learning_rate"]
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max_train_steps = hyperparameters["max_train_steps"]
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output_dir = hyperparameters["output_dir"]
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accelerator = Accelerator(
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gradient_accumulation_steps=gradient_accumulation_steps,
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)
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train_dataloader = create_dataloader(train_batch_size)
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||||
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if hyperparameters["scale_lr"]:
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learning_rate = (
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learning_rate
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||||
* gradient_accumulation_steps
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* train_batch_size
|
||||
* accelerator.num_processes
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||||
)
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||||
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||||
# Initialize the optimizer
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||||
optimizer = torch.optim.AdamW(
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||||
text_encoder.get_input_embeddings().parameters(), # only optimize the embeddings
|
||||
lr=learning_rate,
|
||||
)
|
||||
|
||||
text_encoder, optimizer, train_dataloader = accelerator.prepare(
|
||||
text_encoder, optimizer, train_dataloader
|
||||
)
|
||||
|
||||
# Move vae and unet to device
|
||||
vae.to(accelerator.device)
|
||||
unet.to(accelerator.device)
|
||||
|
||||
# Keep vae and unet in eval model as we don't train these
|
||||
vae.eval()
|
||||
unet.eval()
|
||||
|
||||
# We need to recalculate our total training steps as the size of the training dataloader may have changed.
|
||||
num_update_steps_per_epoch = math.ceil(
|
||||
len(train_dataloader) / gradient_accumulation_steps
|
||||
)
|
||||
num_train_epochs = math.ceil(max_train_steps / num_update_steps_per_epoch)
|
||||
|
||||
# Train!
|
||||
total_batch_size = (
|
||||
train_batch_size
|
||||
* accelerator.num_processes
|
||||
* gradient_accumulation_steps
|
||||
)
|
||||
|
||||
logger.info("***** Running training *****")
|
||||
logger.info(f" Num examples = {len(train_dataset)}")
|
||||
logger.info(f" Instantaneous batch size per device = {train_batch_size}")
|
||||
logger.info(
|
||||
f" Total train batch size (w. parallel, distributed & accumulation) = {total_batch_size}"
|
||||
)
|
||||
logger.info(
|
||||
f" Gradient Accumulation steps = {gradient_accumulation_steps}"
|
||||
)
|
||||
logger.info(f" Total optimization steps = {max_train_steps}")
|
||||
# Only show the progress bar once on each machine.
|
||||
progress_bar = tqdm(
|
||||
range(max_train_steps), disable=not accelerator.is_local_main_process
|
||||
)
|
||||
progress_bar.set_description("Steps")
|
||||
global_step = 0
|
||||
|
||||
for epoch in range(num_train_epochs):
|
||||
text_encoder.train()
|
||||
for step, batch in enumerate(train_dataloader):
|
||||
with accelerator.accumulate(text_encoder):
|
||||
# Convert images to latent space
|
||||
latents = (
|
||||
vae.encode(batch["pixel_values"])
|
||||
.latent_dist.sample()
|
||||
.detach()
|
||||
)
|
||||
latents = latents * 0.18215
|
||||
|
||||
# Sample noise that we'll add to the latents
|
||||
noise = torch.randn(latents.shape).to(latents.device)
|
||||
bsz = latents.shape[0]
|
||||
# Sample a random timestep for each image
|
||||
timesteps = torch.randint(
|
||||
0,
|
||||
noise_scheduler.num_train_timesteps,
|
||||
(bsz,),
|
||||
device=latents.device,
|
||||
).long()
|
||||
|
||||
# Add noise to the latents according to the noise magnitude at each timestep
|
||||
# (this is the forward diffusion process)
|
||||
noisy_latents = noise_scheduler.add_noise(
|
||||
latents, noise, timesteps
|
||||
)
|
||||
|
||||
# Get the text embedding for conditioning
|
||||
encoder_hidden_states = text_encoder(batch["input_ids"])[0]
|
||||
|
||||
# Predict the noise residual
|
||||
noise_pred = unet(
|
||||
noisy_latents, timesteps, encoder_hidden_states
|
||||
).sample
|
||||
|
||||
loss = (
|
||||
F.mse_loss(noise_pred, noise, reduction="none")
|
||||
.mean([1, 2, 3])
|
||||
.mean()
|
||||
)
|
||||
accelerator.backward(loss)
|
||||
|
||||
# Zero out the gradients for all token embeddings except the newly added
|
||||
# embeddings for the concept, as we only want to optimize the concept embeddings
|
||||
if accelerator.num_processes > 1:
|
||||
grads = (
|
||||
text_encoder.module.get_input_embeddings().weight.grad
|
||||
)
|
||||
else:
|
||||
grads = text_encoder.get_input_embeddings().weight.grad
|
||||
# Get the index for tokens that we want to zero the grads for
|
||||
index_grads_to_zero = (
|
||||
torch.arange(len(tokenizer)) != placeholder_token_id
|
||||
)
|
||||
grads.data[index_grads_to_zero, :] = grads.data[
|
||||
index_grads_to_zero, :
|
||||
].fill_(0)
|
||||
|
||||
optimizer.step()
|
||||
optimizer.zero_grad()
|
||||
|
||||
# Checks if the accelerator has performed an optimization step behind the scenes
|
||||
if accelerator.sync_gradients:
|
||||
progress_bar.update(1)
|
||||
global_step += 1
|
||||
|
||||
logs = {"loss": loss.detach().item()}
|
||||
progress_bar.set_postfix(**logs)
|
||||
|
||||
if global_step >= max_train_steps:
|
||||
break
|
||||
|
||||
accelerator.wait_for_everyone()
|
||||
|
||||
# Create the pipeline using using the trained modules and save it.
|
||||
if accelerator.is_main_process:
|
||||
pipeline = StableDiffusionPipeline(
|
||||
text_encoder=accelerator.unwrap_model(text_encoder),
|
||||
vae=vae,
|
||||
unet=unet,
|
||||
tokenizer=tokenizer,
|
||||
scheduler=PNDMScheduler(
|
||||
beta_start=0.00085,
|
||||
beta_end=0.012,
|
||||
beta_schedule="scaled_linear",
|
||||
skip_prk_steps=True,
|
||||
),
|
||||
safety_checker=StableDiffusionSafetyChecker.from_pretrained(
|
||||
"CompVis/stable-diffusion-safety-checker"
|
||||
),
|
||||
feature_extractor=CLIPFeatureExtractor.from_pretrained(
|
||||
"openai/clip-vit-base-patch32"
|
||||
),
|
||||
)
|
||||
pipeline.save_pretrained(output_dir)
|
||||
# Also save the newly trained embeddings
|
||||
learned_embeds = (
|
||||
accelerator.unwrap_model(text_encoder)
|
||||
.get_input_embeddings()
|
||||
.weight[placeholder_token_id]
|
||||
)
|
||||
learned_embeds_dict = {
|
||||
args.placeholder_token: learned_embeds.detach().cpu()
|
||||
}
|
||||
torch.save(
|
||||
learned_embeds_dict, os.path.join(output_dir, "learned_embeds.bin")
|
||||
)
|
||||
|
||||
|
||||
import accelerate
|
||||
|
||||
accelerate.notebook_launcher(
|
||||
training_function, args=(text_encoder, vae, unet), num_processes=1
|
||||
)
|
||||
|
||||
# Set up the pipeline
|
||||
pipe = StableDiffusionPipeline.from_pretrained(
|
||||
hyperparameters["output_dir"],
|
||||
# torch_dtype=torch.float16,
|
||||
)
|
||||
|
||||
all_images = []
|
||||
for _ in range(args.num_inference_samples):
|
||||
images = pipe(
|
||||
[args.prompt],
|
||||
num_inference_steps=args.inference_steps,
|
||||
guidance_scale=7.5,
|
||||
).images
|
||||
all_images.extend(images)
|
||||
|
||||
# output_path = os.path.abspath(os.path.join(os.getcwd(), args.output_dir))
|
||||
if not os.path.isdir(args.output_dir):
|
||||
os.mkdir(args.output_dir)
|
||||
|
||||
[
|
||||
image.save(f"{args.output_dir}/{i}.jpeg")
|
||||
for i, image in enumerate(all_images)
|
||||
]
|
||||
Reference in New Issue
Block a user