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release_0.
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README.md
39
README.md
@@ -3,21 +3,7 @@
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This is a fork of CompVis/stable-diffusion, the wonderful open source
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text-to-image generator.
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The original has been modified in several minor ways:
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## Simplified API for text to image generation
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There is now a simplified API for text to image generation, which
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lets you create images from a prompt in just three lines of code:
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~~~~
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from ldm.simplet2i import T2I
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model = T2I()
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outputs = model.text2image("a unicorn in manhattan")
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~~~~
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Outputs is a list of lists in the format [[filename1,seed1],[filename2,seed2]...]
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Please see ldm/simplet2i.py for more information.
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The original has been modified in several ways:
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## Interactive command-line interface similar to the Discord bot
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@@ -63,8 +49,25 @@ weights (512x512) and the older (256x256) latent diffusion weights
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identical to those used in the Discord bot, except you don't need to
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type "!dream". Pass "-h" (or "--help") to list the arguments.
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This script also provides an img2img feature that lets you seed your
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creations with a drawing or photo.
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For command-line help, type -h (or --help) at the dream> prompt.
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## Simplified API for text to image generation
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There is now a simplified API for text to image generation, which
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lets you create images from a prompt in just three lines of code:
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~~~~
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from ldm.simplet2i import T2I
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model = T2I()
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outputs = model.text2image("a unicorn in manhattan")
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~~~~
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Outputs is a list of lists in the format [[filename1,seed1],[filename2,seed2]...]
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Please see ldm/simplet2i.py for more information.
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## Workaround for machines with limited internet connectivity
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My development machine is a GPU node in a high-performance compute
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@@ -116,14 +119,16 @@ to send me an email if you use and like the script.
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# Original README from CompViz/stable-diffusion
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*Stable Diffusion was made possible thanks to a collaboration with [Stability AI](https://stability.ai/) and [Runway](https://runwayml.com/) and builds upon our previous work:*
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[**High-Resolution Image Synthesis with Latent Diffusion Models**](https://arxiv.org/abs/2112.10752)<br/>
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[**High-Resolution Image Synthesis with Latent Diffusion Models**](https://ommer-lab.com/research/latent-diffusion-models/)<br/>
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[Robin Rombach](https://github.com/rromb)\*,
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[Andreas Blattmann](https://github.com/ablattmann)\*,
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[Dominik Lorenz](https://github.com/qp-qp)\,
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[Patrick Esser](https://github.com/pesser),
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[Björn Ommer](https://hci.iwr.uni-heidelberg.de/Staff/bommer)<br/>
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which is available on [GitHub](https://github.com/CompVis/latent-diffusion).
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**CVPR '22 Oral**
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which is available on [GitHub](https://github.com/CompVis/latent-diffusion). PDF at [arXiv](https://arxiv.org/abs/2112.10752). Please also visit our [Project page](https://ommer-lab.com/research/latent-diffusion-models/).
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[Stable Diffusion](#stable-diffusion-v1) is a latent text-to-image diffusion
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@@ -197,7 +197,7 @@ class T2I:
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shape = [self.latent_channels, height // self.downsampling_factor, width // self.downsampling_factor]
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samples_ddim, _ = sampler.sample(S=steps,
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conditioning=c,
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batch_size_size=batch_size,
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batch_size=batch_size,
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shape=shape,
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verbose=False,
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unconditional_guidance_scale=cfg_scale,
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